ASP08

Analog Signal Processing
Dr.-Ing. Tuan Do-Hong

Department of Telecommunications Engineering
Ho Chi Minh City University of Technology
do-hong@hcmut.edu.vn
September, 2008
Course Overview
0.1
Course Description
Introduction to analog signal processing, with an emphasis on underlying concepts from circuit and system analysis: linear systems, review of
elementary circuit analysis, dierential equation models of linear circuits
and systems, Laplace transform, convolution, stability, phasors, frequency
response, Fourier series, Fourier transform, active lters and AM radio.
0.2
Prerequisite
Calculus. Dierential equations. Physics-bases treatment of electricity and

magnetism. Introductory exposure to circuit analysis.
0.3
Goals
To introduce fundamentals of analog signal processing, with major emphasis on circuit analysis, dierential equations, convolutions, Fourier methods,
and applications in ltering and AM radio.
0.4
Topics
Examples of signals and signal processing systems

Analog linear shift-invariant systems
Circuits and linear systems
Review of DC circuit analysis: KCL, KVL, dependent sources
Capacitors and inductors as circuit elements
Op-amp circuits
2
0.5. ASSESSMENT
Characterization and solution of LSI systems via linear, constantcoecient dierential equations
Transient response of RLC circuits
Complex numbers and functions of a complex variable
One-sided Laplace transform
Impedance
Laplace transform solution of dierential equations
General form of solution to a dierential equation
Transfer function and block diagrams
Impulse as a generalized function
Convolution
Stability
Sinusoidal steady-state analysis and phasors
Steady-state analysis of circuits
Frequency response
Fourier series
Fourier transform
Design of active lters
AM radio
Sampling theorem and overview of digital signal processing
0.5
Assessment
Home assignments: 20%

Mid-term exam: 20%
Final exam: 60%
0.6
Texts
B. P. Lathi, Signal Processing and Linear Systems, Berkeley-Cambridge

Press, 1998.
David E. Johnson, John L. Hilburn, Johnny R. Johnson, Basic Electric
Circuit Analysis, 3rd edition, Prentice-Hall, 1986.
0.7
Course Content
1. Fundamentals of Circuit Analysis

Circuit introductory.
Circuit reduction techniques.
Circuit analysis techniques.
Real sources, controlled sources and maximum power transfer.
Operational ampliers.
Time-dependent circuits.
First-order circuits.
Second-order circuits.
2. Introduction to Signals and Systems
Classication of signals.
Signal operations and models.
Classication of systems.
System model: Input-output Description.
3. Time-Domain Analysis of Continuous-Time Systems
System response to internal conditions: Zero-input response.
Unit impulse response h(t).
System response to external conditions: Zero-state response.
Classical solution of dierential equation.
System stability.
4. Signal Representation by Fourier Series
0.7. COURSE CONTENT
Signals and vectors.

Signal representation by orthogonal signal set.
Trigonometric Fourier series.
Exponential Fourier series.
LTIC system response to periodic inputs.
5. Continuous-Time Signal Analysis: Fourier Transform
Aperiodic signal representation by Fourier integral.
Properties of Fourier transform.
Signal transmission through LTIC systems.
Ideal and practical lters.
Application to communications: Amplitude Modulation (AM).
6. Continuous-Time Signal Analysis Using Laplace Transform
Laplace transform and properties.
Solution of dierential and integro-dierential equations.
Analysis of electrical networks: Transformed network.
Block diagram.
Application to feedback and controls.
7. Frequency Response and Analog Filters
Frequency response of an LTIC system.
Bode plots.
Filter design by placement of poles and zeros of H(s).
Butterworth lters.
Chebyshev lters.
8. Sampling
Sampling theorem.
Numerical computation of Fourier transform: Discrete Fourier
Transform (DFT).
Fast Fourier Transform (FFT).
Chapter 1
Fundamentals of Circuit Analysis
1.1
1.1.1
Circuit Introductory
Circuit Variables
Scientic investigation of static electricity was done in late 1700s and

Coulomb is credited with most of the discoveries. He found that electric charges have two attributes: amount and polarity. There are two type
of charges opposite charges attract and similar polarity ones repel each
other. Charge polarity is indicated by positive and negative signs because
positive and negative charges cancel each other when brought together. As
a results, the electric charge can be described by an algebraic number, q,
with units of Coulomb (C).
Because opposite charges attract each other, energy is expanded to separate them from each other. This energy is stored in the electric eld between
the two reservoir of separated charges and is recovered when the charges
are allowed to come together. The stored energy per unit charge is called
the voltage or potential dierence between the two reservoir of charges:
v=
dW
dq
(1.1)
where the unit is Volt (V = CJ ).

Note that we need two reservoir of charges. So voltage is between two
points. We also use + and - signs to indicate the direction for measuring
v. From denition of voltage above, w is energy needed to move a positive
charge from - reservoir to + reservoir.
One can dene a reference point for measuring voltages (typically shown
6
1.1. CIRCUIT INTRODUCTORY
as ground). The voltage between any point and this reference point is call
the potential of that point. It is always assumed that + is at the point and
the - sign is at the reference point. Therefore, there is no need to indicate
+ and - signs for potential. Voltage between two points is the dierence
between the potential of the two points (see Figure 1.1).
Voltage between two charge reservoirs is analogous to height dierence
between two uid reservoir and the same way, the potential of each point
is analogous to its elevation compared to some reference (e.g., sea level).
Potential
V1
Elevation
+
V2
v = v12 = V1 - V2
v = v21 = V2 - V1
v = -v
h2
h1
Sea level
h = h1 h2
Figure 1.1: Voltage between two charge reservoirs.
If we connect the charge reservoirs, electric charges ow from one to

other. The rate of the charge ow through a specic area is called the
electric current:
dq
(1.2)
i=
dt
with the current owing in the direction of the charge ow (it means that
a positive current is associated with the ow of positive charge). The unit
of current is Ampere (A = C/s).
In principle, electric charges generate an electric eld and motion of the
charged particles (current) generates a magnetic eld. This electromagnetic
eld interacts with all charges and aect them. The behavior of such a
system is described by Maxwells equation. Solution of Maxwells equations,
CHAPTER 1. FUNDAMENTALS OF CIRCUIT ANALYSIS
however, is dicult and not needed expect for some cases (propagation of
electromagnetic wave and light, antennas, etc.).
In most relevant engineering cases the problem can be greatly simplied
by noting that electric charges preferentially ow through a conductor (or
a semiconductor) as opposed to vacuum, air, or any insulator. In this case,
the system can be described as a circuit containing circuit elements and
connecting ideal wires. Circuit theory is the scientic discipline that describes behavior of circuits be and is built upon the following assumptions:
All of the electromagnetic phenomena occurs inside each circuit element. They communicate with the outside world only through the
voltage across and current that particular element.
Circuit elements are connected to each other with ideal wires that do
not impede ow of charge. They can be stretched (making them longer
or shorter, for example) without any eect on the circuit.
Net amount of charge cannot be accumulated in any circuit element
or any location in the circuit. If a net charge of q enters a circuit
element, the same amount of charge should leave the element. This
means: (a) a circuit element should have at least two terminals, (b)
Because current travels through the system at a good fraction of speed
of light, we can safely assume that the total current entering a circuit
element is exactly equal to the current leaving that element at any
instant time.
i
+
1
Circuit element
v
2
Ideal wire
i
Figure 1.2: Circuit element
with two terminals.
In the context of circuit denition above, the important physical quantities are current and voltages (and electric power). These are the circuit
variables that form the basis for communication between circuit elements.
The value of voltage between two points, v, is incomplete unless the +
and - signs are assigned to the two points. Similarly, the value of current
is incomplete unless a reference current direction is assigned. However, as
v and i are algebraic numbers, the assigned position of + and - for voltage
and direction for the current is arbitrary. The sign of v and i ips according
to this choice.
Internal of each circuit element impose a relationship between the current owing in the element and the voltage across that element. This relationship is called the element law or i-v characteristics. While the choice
of reference directions for current through and voltage across an element
is arbitrary, the element law or i-v characteristics of an element depends
on the reference directions. To see this point, note that there exists four
possible choices for voltage and current directions in a two-terminal circuit
element:
i
-
(a )
v
2
+
1
( b)
Figure 1.3: (a) Passive sign convention. (b) Active sign convention.
In two left cases, the current direction is marked such that it ows from
+ to - signs. This case is called the passive sign convention. In the two
right cases, current ows from - to + signs. This case is referred to as active
sign convention. Obviously, the i-v characteristic will be dierent in each
case as the sign of current or voltage is switched. To resolve this confusion,
we follow this convention:
10
i-v characteristics of two-terminal element are written assuming passive sign convention.
Use passive sign convention when marking current and voltages in a
circuit as much as possible.
So, in marking references for voltages and currents in the circuit, it is
best to arbitrarily choose either the voltage or current reference direction
and then use passive sign convention to mark the other variable. Note that
it may not be possible or practical to mark every element using passive
sign convention. In this case, a good rule is to mark every element except
voltage and current sources using passive sign convention.
Electric power produced or absorbed in an element
The power, denoted by P , is dened as
P =
dW
dt
(1.3)
where the unit is Watt (W = J/s). Since v = dW/dq and i = dq/dt (both
total derivatives), then:
P =
dW
dq
dW
=
=iv
dt
dq
dt
(1.4)
The quantities v and i are generally function of time. Therefore P is a

time-varying quantity, sometimes called the instantaneous power.
Using the denition of the voltage, one nds that if we use passive sign
convention:
P < 0: Element is producing power.
P > 0: Element is absorbing power.
Obviously, if we use active sign convention:
P < 0: Element is absorbing power.
P > 0: Element is producing power.
11
i
A
+
Figure 1.4: Absorbed
and supplied powers.
Example
Find the power absorbed or supplied by elements A and B if i = 25A and
v = 120V.
For element A,
P = vi = 25 120 = 3, 000W = 3KW
Since element A reference directions follow passive sign convention and
P > 0, element A is absorbing 3 KW of power.
For element B,
P = vi = 25 120 = 3, 000W = 3KW
Since element B reference directions follow active sign convention and P > 0,
element B is supplying 3 KW of power.
1.1.2
Kirchho Laws
A circuit consists of circuit elements attached to each other with ideal

wires or connectors.
Node: A node is a point in the circuit that is connected to two or more
circuit elements with ideal wires.
Loop: A loop is a closed path in a circuit through at least two circuit
elements which return to starting node without passing through any node
twice.
12
Kirchho Current Law (KCL): KCL follows from our assumption

that no net charge can be accumulated at any point in the circuit (or in a
node): Sum of currents entering a node is equal to sum of currents leaving
a node.
Example: Write KCL for the marked node.
+i1 i2 i3 + i4 = 0
i1
i2
i3
i4
Figure 1.5: Example for KCL.
Kirchho Voltage Law (KVL): KVL follows from the denition of

the voltage between points. Voltage is dened as the amount energy to
move unit charge from one point to another. Zero net energy should be
expanded if a charge q is moved around a closed loop and is returned back
to its original position, i.e.,

W = q
v = 0
v = 0
(1.5)
Algebraic sum of all voltages around any closed path in a circuit is zero.
Example: Write KVL for the marked loop.
+vb + vc va = 0
13
vb
va
vc
Figure 1.6: Example for KVL.
KVLs and KCLs are constraints on circuit variables which arise because
of the circuit arrangement (attachment of connection wires). In addition,
internal of each circuit element impose a relationship between the current
owing in the element and the voltage across that element (element laws
or i-v characteristics). Combination of these two set of constraints results
in a unique set of values for the circuit variables (currents and voltages).
In a circuit with E elements, there are 2E circuit variables (i and v for
each element). We need 2E equations to nd these circuit variables. If the
circuit has N nodes, we can write (N 1) KCL equations (KCL on the
N th node is exactly the sum of KCL on the other (N 1) nodes). We can
also write E N + 1 independent KVLs and E i-v characteristic equations.
Totally, it is 2E equations.
If the i-v characteristic is a linear relationship between i and v, the
element is a linear element. If a circuit is made of linear element, the
resulting set of 2E equations in 2E variables for a linear algebraic set of
equations. In this course, we only use linear elements.
1.1.3
Linear Circuit Elements
Resistor
i-v characteristic: v = Ri.
Resistance: R. Unit: Ohm ().
Conductance: G = 1/R. Unit: Siemens (S).
P = vi = Ri2
(1.6)
14
P > 0: Resistor always absorbs power.

i
i = v/ R
+ v -
R
Figure 1.7: Resistor: Circuit symbol and i-v characteristic.
Independent Voltage Source (IVS)

i-v characteristic: v = vs for any i.
+ i
v
-
+
-
v = vs
vs
vs
Figure 1.8: Independent voltage source: Circuit symbol and

i-v characteristic.
Independent Current Source (ICS)

i-v characteristic: i = is for any v.
Short Circuit
i-v characteristic: v = 0 for any i.
Note that a short-circuit element is a special case of a resistor (with
R = 0) or an ideal voltage source (with vs = 0).
Open Circuit
i-v characteristic: i = 0 for any v.
Note that a open-circuit element is a special case of a resistor (with
R ) or an ideal current source (with is = 0).
15
+ i
v
is
is
i = is
v
Figure 1.9: Independent current source: Circuit symbol and

i-v characteristic.
+ i
v
v
Figure 1.10: Short circuit: Circuit symbol and i-v characteristic.
+ i
v
-
Figure 1.11: Open circuit: Circuit symbol and i-v characteristic.
16
Switch
+ i
v
-
i-v characteristic:
Switch open: Open circuit i = 0 for any v.
Switch close: Short circuit v = 0 for any i.
1.1.4
Example
Consider the circuit as in Figure 1.12. Find i1
25
+ v0 25
i0
i2
10
i1
+
70
180V
180V
i5
i0
10
+ v2 - i 1
+
v1
-
+ v3 70
i3
i4
+
v4
-
(b)
(a)
Figure 1.12: Example of circuit analysis using KVL and KCL. (a) Original circuit. (b) Circuit
with nodes and loops.
We rst mark the circuit variables consisting of 11 unknowns:

v0 , v1 , v2 , v3 , v4 , i0 , i1 , i2 , i3 , i4 , i5
17
identify the nodes, and write the i-v characteristics equations:

v0
v1
v2
v3
v4
=
=
=
=
=
25i0
70i1
5i2
10i3
8i4
Note that we have 6 elements so we have 12 circuit variables. However,

the IVS element law species the voltage across its terminal, so we have 11
circuit variables which have to be found.
The circuit has four nodes. Since we need to write KCL only in N 1
nodes, we choose not write KCL in the bottom node. The KCLs are:
i0 + i2 + i5 = 0
i3 + i1 i2 = 0
i4 i3 i0 = 0
We need to write NKV L = 11 5 3 = 3 KVLs. Choosing 3 loops with
smallest number elements, we get:
= 0
+v0 v3 v2
+v2 + v1 (180) = 0
= 0
+v3 + v4 v1
Above are eleven equations in eleven unknowns that can be solved. The
numbers of equations to be solved can be halved by using i-v characteristics
equations to substitute in KVLs and KCLs. For example, if we substitute
for voltages from i-v characteristics equations in KVLs, we get 6 equations
in 6 unknown currents:
i0 + i2 + i5
i3 + i1 i2
i4 i3 + i0
25i0 10i3 5i2
5i2 + 70i1
10i3 + 8i4 70i1
=
=
=
=
=
=
0
0
0
0
180
0
Note: As can be seen, even relatively small circuits results in a large

number equations to be solved. Several techniques for reducing the number
of equation to be solved are introduced next.
18
1.2
Circuit Reduction Techniques
Combination of KVLs, KCLs, and i-v characteristics equations result in

a set of linear equations for the circuit variables. While the above set of
equation is complete and contains all necessary information, even small circuits require a large number of simultaneous equations to be solved as seen
previously. We learned that we can use i-v characteristics equations when
we are marking the circuit variables and reduce the number of equations
to be solved to only number of KCLs and KVLs. We will learn later two
methods, node-voltage and mesh-current, which reduce the number of equations to be solved further to either number of KCLs or number of KVLs.
This is the best we can do in this direction.
As it is easier to solve smaller sets of equations (e.g., it is easier to
solve two sets of two equations in two unknown as compared to a set of 4
equations in 4 unknowns), one can break up the circuit into smaller pieces
and solve each individually and assemble back the whole circuit. We can
use this principle to combine circuit elements and make a much smaller
circuit. These techniques are described below.
Recall that we are using lumped circuit elements, i.e., circuit elements
communicate to outside world and other circuit elements only through i
and v. Conversely, the outside world (the rest of the circuit) communicate
with the circuit element through i and v. This means, for example, that
a resistor in a circuit is viewed by the rest of the circuit as a black box
with an i-v characteristics of v = Ri. The rest of the circuit does not know
some elements what is inside the box. In fact, we can replace the resistor
with any black box (containing whatever) with the same i-v characteristics
of v = Ri and the rest of the circuit behaves exactly the same.
Alternatively, if a black box containing many circuit elements is attached
to a circuit and has an i-v characteristics of v = 5i, we can replace this black
box with a 5 resistor with no change in the circuit behavior.
This observation allows the circuit to be divided into two or many parts
and each solved independently. We dene a box containing several element
a subcircuit or a device. The above gure shows a two-terminal device or
subcircuit.
1.2. CIRCUIT REDUCTION TECHNIQUES
19
i
+
Rest of
the circuit
Subcircuit
A box with some

circuit elements
i
+
Rest of
the circuit
v
Equivalent
subcircuit
Figure 1.13: Subcircuit and rest of the circuit.
20
Subcircuits play an important role in linear circuit theory. Thevenin

theorem states that any subcircuit containing linear circuit element has an
i-v characteristics of Av +Bi = C (where A, B, and C are constants) and it
can be reduced to a subcircuit containing at most two linear circuit elements
(Thevenin and Norton Forms). We will discuss Thevenin Theorem later.
Below, we explore subcircuits in the context of elements that are in series
or in parallel. In each case, we nd the i-v characteristics of the subcircuit
and use that to nd the equivalent element.
Two Resistors in Series
KCL : i = i1 = i2
KVL : v + v1 + v2 = 0
i-v :
v1 = R1 i1 = R1 i
v2 = R2 i2 = R2 i
i
+
v
-
R1
i1
+ v1 - +
v2
-
R2
i2
Figure 1.14: Two resistors

in series.
Substituting from i-v characteristics equations in KVL, we get

v
= R1 i + R2 i
= (R1 + R2 )i
= Req i
Req = R1 + R2
So, a subcircuit containing two resistors in series has an i-v characteristics of the form v = Req i and is equivalent to a resistor, Req = R1 + R2 .
The above can be easily extended: k resistors in series are equivalent to
one resistor with

Req =
k
21
Rj
(1.7)
j=1
A Resistor in Series with a Current Source

KCL : i = i1 = is
KVL : v + v1 + v2 = 0
v = R1 is = v2
i
+
v
-
i1
+ v1 - +
v2
-
is
Figure 1.15: A resistors in

series with a current source.
The i-v characteristics of ICS states that its current is is, independent
of its voltage (v2 ). Above equations show that the i-v characteristics of
subcircuit is i = is and is independent of voltage v = v1 + v2 (as value of
v1 can be anything). The equivalent subcircuit is an independent current
source with strength is .
i
+
R
+
vs
v
-
Figure 1.16: Thevenin form.
A Resistor in Series with a Voltage Source

This is the Thevenin form and cannot be reduced further.
22
A Voltage Source in Series with a Voltage Source

KVL : v + vs1 + vs2 = 0
v = vs1 + vs2
vs1
+
vs2
v
-
Figure 1.17:
A voltage
source in series with a voltage source.
The i-v characteristics of IVS states that their voltages are vs1 and vs2 ,
respectively, independent of their currents. Above equation shows that the
i-v characteristics of subcircuit is v = vs1 +vs2 and is independent of current
i. The equivalent subcircuit is an independent voltage source with strength
vs1 + vs2 .
A Voltage Source in Series with a Current Source
is 1
+ v1 -
+
vs2
v
-
i2
Figure 1.18:
A voltage
source in series with a current source.
23
KCL : i = is1 = is2

KCL : v + v1 + vs2 = 0
v = v1 + vs2
The i-v characteristics of ICS states that its current is is1 , independent
of its (v1 ). Above equations show that the i-v characteristics of subcircuit
is i = is1 and is independent of voltage v = v1 + vs2 (as value of v1 can be
anything). The equivalent subcircuit is an independent current source with
strength is1 .
Note: A current source in series with any element reduced to a current
source. Series element all have the same current and the current source
requires the current through to be equal to its strength.
A Current Source in Series with a Current Source
KCL : i = is1 = is2
is 1
+
v
is 2
Figure 1.19: Two current

sources in series.
A current source in series with any element reduces to a current source.

However, in the case of two current sources in series, KCL requires is1 = is2 .
Thus, two current sources can be attached in series only if is1 = is2 . If so,
the equivalent subcircuit is an independent current source with strength
is = is1 = is2 .
24
Two Resistors in Parallel

KCL : i + i1 + i2 = 0
KVL : v = v1 = vs2
i-v :
v1 = R1 i1
v2 = R2 i2
Substituting from i-v characteristics equations in KCL, and using v =
v1 = v2 , we get
1
1
1
v
=
+
i=
Req
Req
R1 R2
So, a subcircuit containing two resistors in parallel has an i-v characteristics of the form v = Req i and is equivalent to a resistor, 1/Req =
1/R1 + 1/R2 .
The above can be easily extended: k resistors in series are equivalent to
one resistor with
k

1
1
=
(1.8)
Req
R
j
j=1
A Resistor in Parallel with a Current Source
i
+
v
-
+
v1
-
i1
R1
+
v2
-
is
Figure 1.20: Norton form.
This is the Norton form and cannot be reduced further. We will show
later that Norton and Thevenin forms are equivalent.
A Resistor in Parallel with a Voltage Source
KCL : i + i1 + i2 = 0
KVL : v = v1 = vs
25
i
+
v
-
+
v1
-
i1 i2
R1
-
vs
Figure 1.21:
A resistor
in parallel with a voltage
source.
The i-v characteristics of IVS states that its voltage is vs , independent

of its current, i2 . Above equations show that the i-v characteristics of
subcircuit is v = vs independent of current i. The equivalent subcircuit is
an independent voltage source with strength vs .
A Current Source in Parallel with a Current Source
KCL : i + is1 + is2 = 0
i = is1 + is2
i
+
v
v1
-
is 1
+
v2
-
is 2
Figure 1.22: Two current

sources in parallel.
The i-v characteristics of ICSs state that their currents are is1 and is2 ,
respectively, independent of their voltage, v. Above equations show that
the i-v characteristics of subcircuit is i = is1 + is2 independent of value of
voltage v. The equivalent subcircuit is an independent voltage source with
strength is = is1 + is2 .
26
A Current Source in Parallel with a Voltage Source

KCL : i + i1 + i2 = 0
KVL : v = vs = v2
i
i1
+
v
-
vs
-
+
v2
-
is
Figure 1.23:
A current
source in parallel with a
voltage source.
The i-v characteristics of IVS states that its voltage is vs , independent

of its current, i1 . Above equations show that the i-v characteristics of
subcircuit is v = vs and is independent of current i = i1 + is (as value of
i1 can be anything). The equivalent subcircuit is an independent voltage
source with strength vs .
Note: A voltage source in parallel with any element reduces to a voltage
source. Parallel elements all have the same voltage and the voltage source
requires its voltage across to be equal to its strength.
A Voltage Source in Parallel with a Voltage Source
KVL : v = vs1 = vs2
A voltage source in parallel with any element reduces to a voltage source.
However, in the case of two voltage sources in parallel, KVL requires vs1 =
vs2 . Thus, two voltage sources can be attached in parallel only if vs1 = vs2 .
If so, the equivalent subcircuit is an independent voltage source with strength
vs = vs1 = vs2 . circuit conguration arises because we are dealing wit
idealized circuit element. We will discuss real sources later.
27
i
i1
+
+
vs1
-
i2
+
vs2
-
Figure 1.24: Two voltage

sources in parallel.
Thevenin and Norton Forms and Source Transformation

RT
+
RN
vT
iN
Figure 1.25: Thevenin and Norton forms are equivalent if

RT = RN and vT = iN RT .
Thevenin and Norton forms are equivalent if RT = RN and vT = iN RT .

One can replace one with the other. This is called source transformation
and is helpful in rearranging other elements in the circuit and sometime
arriving at more element being in series and parallel.
Voltage Divider
The two resistors can be replaced by an equivalent resistor, Req = R1 + R2 .
Thus
vs
vs = Req i i =
Req
v1 = iR1 =
R1
vs
Req
v2 = iR2 =
R2
vs
Req
28
i
R1
+
v1
R2
+
v2
+
vs
-
Figure 1.26:
vider.
Voltage di-
Also,
R1
v1
=
v2
R2
This circuit is called a voltage divider as the two resistors divide the
voltage of the IVS between them proportional to their values. This circuit
can be extended by adding more resistors to the circuit and get more reference voltages. This circuit is used extensively in electronic circuits. The
basic reason is that power supplies are bulky and/or expensive. Typically,
one power supply with one voltage is provided. On the other hand, more
than one voltage may be needed for the circuit to operate properly.
Example: A battery operated radio has a 9V battery. Part of radio
circuits require a 6V supply. Design a voltage divider circuit to supply 6V
voltage to these circuits. The desired circuit is the voltage divider circuit
above with vs = 9V and v2 = 6V. Then,
v2 =
R2
R2
R2
6
vs 6 =
9
=
Req
R1 + R2
R1 + R2
9
This is one equation in two unknowns and one is free to choose one
parameter. For example, choosing R1 = 1K, we get R2 = 2K. Voltage
dividers are aected by the load current drawn from them (see gure). The
voltage divider formula can only be used for the circuit if il i.
29
i
R1
+
v1
R2
- iL
+
v2
+
vs
-
Load
Figure 1.27: Voltage divider with load.
Current Divider
The two resistors can be replaced by an equivalent resistor,
1/Req = 1/R1 + 1/R2
is
v
-
+
v
-
i1 i2
R2
R1
+
v
-
Figure 1.28: Current divider.
Thus
v = Req i
v = i1 R1 i1 =
v
Req
=
is
R1
R1
v = i2 R2 i2 =
v
Req
=
is
R2
R2
Also,
R2
i1
=
i2
R1
30
This circuit is called a current divider as the two resistors divide the
current of the ICS between them (inversely proportional to their values).
This circuit can be extended by adding more resistors to the circuit and
get more reference currents.
Wheatstone Bridge
R1
+
vs
-
A
R2
RA
+ vm + +
v2 vB
- -
B
RB
Figure 1.29: Wheatstone bridge.
A typical Ohm-meter measures the resistance of a resistor by using the

Ohms Law. It applies a known voltage of vs across the resistor, measures
the current owing through the resistor, and its dial are set to convert
the measured value of current into the value of resistance by using R =
vs /imeasured . (This is why one cannot measure the value of a resistor while
it is attached in a circuit, Ohm-meter works only if the resistor is not
attached to anything but the meter.)
A typical digital multi-meter measure resistance within an accuracy of
about 1%. In some cases, higher accuracy is needed. Resistor bridges are
used for this purpose and they are made of two voltage divider circuits put
in parallel with each other. The bridges operate based on the fact that
while it is dicult to measure the dierence between 1 and 1.01 V or 1
and 1.01 A distinctly (they are only 1% apart and within the accuracy of
meter), it is easy to measure 0.01 V. A most widely used bridge is the
Wheatstone bridge shown. It consists of two voltage divider circuits with a
voltmeter measuring the voltage between points A and B (denoted by vm ).
1.3. CIRCUIT ANALYSIS TECHNIQUES
31
We note from voltage divider formulas

v2 =
R2
vs
R1 + R2
vB =
RB
vs
RA + RB
KVL : v2 + vm + vB = 0 vm = v2 vB = vs
RB
R2
R1 + R2 RA + RB
Measuring resistance accurately: Suppose the unknown resistance is RB .

Two known and accurate resistors R1 and RA are chosen. An accurate but
variable resistor is used for R2 . The variable resistor R2 is varied until the
meter read zero voltage for vm . Then

RB
R2
R2
= 0 RB = R A
vm = vs
R1 + R2 RA + RB
R1
Note that we do not need to know value of vs to nd RB . This way, RB
is measure within the accuracy of resistors, R1 , R2 , and RA .
1.3
Circuit Analysis Techniques
In this section, we learn two methods, node-voltage and mesh-current methods, which reduce the number of equations to either number of KCLs or
number of KVLs.
1.3.1
Node-Voltage Method
The node-voltage method is based on following idea. Instead of solving for

circuit variables, i and v of each element, we solve for a dierent set of
parameters, node voltages in this case, which automatically satisfy KVLs.
As such, we do not need to write KVLs and only need to solve KCLs.
Recall denition of potential. Potential was dened as the voltage between any point and the reference point (usually called ground or common).
Take a circuit and identify the nodes on the circuit. Dene the node voltage
as the potential of that node (voltage between that node and the reference).
In the circuit shown, the node voltages are denoted as v1 , v2 , and v3 . The
voltage across each element is simply the voltage of the node attached to
32
v1
vb
v2
va
vc
v3
Reference node
( g r o u nd )
Figure 1.30: Node voltage: voltage between that node and the reference.
+ terminal of the element minus

minal of the element:
va
vb
vc
the voltage of the node attached to - ter= v1 v2

= v1 v3
= v2 v3
Note that using node voltage denitions, KVLs will be automatically

satised:
KVL : vb + va + vc = 0
(v1 v3 ) + (v1 v2 ) + (v2 v3 ) = 0 0 = 0
As the choice of the reference point (ground) is arbitrary, we can choose
the reference point to be any node on the circuit. For example, we could
choose node 3 in the above circuit as the reference mode (v3 = 0). The
number of node voltages is reduced to 2. The voltages across the three
elements are
va = v1 v2
vb = v1
vc = v2
and KVL is still automatically satised. With this simplication, we have
N 1 node voltages in each circuit where N is the number of nodes. Recall
33
that number of independent KCLs is N 1. If we can write the current

in each element in terms of the node voltages, then the circuit reduced to
N 1 KCLs in N 1 node voltages.
i1
v1
+
vs
-
is
i2
v2
i1
v1
+ v1 2 R
- v2 1 +
i2
v2
Figure 1.31: Node voltages of current source and resistor.
Current Source: The current owing through current source is independent of voltage across it and, therefore, is independent of node voltages.
From the gure, it is obvious that the current leaving node 1 is i1 = is and
current leaving node 2 is i2 = is .
Voltage Source: The voltage across a voltage source is constant and
is independent of current owing through it. The i-v characteristics of
the voltage source cannot be used to nd the current owing through it.
Voltage sources should be treated in a special manner that will be discussed
later.
Resistor: Resistors follow Ohms Law. We need to calculate the voltage
across the resistor rst and nd the current using the Ohms Law. For
example, if we want to nd current leaving node 1, i1 , we rst nd voltage
v12 :
v12
v1 v2
v12 = v1 v2 i1 =
=
R
R
Alternatively, if we want to nd the current leaving node 2, i2 , we calculate
v21 rst:
v21
v2 v1
v21 = v2 v1 i2 =
=
R
R
In summary, we have found that we can dene node voltages as circuit
variables and:
Voltage across each element can be written in terms of node voltages.
34
Current through each element (resistor and ICS) can be written in

terms of node voltages using the i-v characteristics of that element.
KVLs are automatically satised.
By using node voltages, we have reduced the number of variables to
N 1 node voltages (one node is the reference node with zero voltage)
and N 1 KCLs to solve.
v1
is1
i2
R2
i2
i1
i3
R1
R3
v2
is2
Figure 1.32: Example for node-voltage method.
Example: The circuit in Figure 1.32 has three nodes. We choose the
node at the bottom as the reference (ground).The voltages at the other
nodes are denoted by v1 and v2 . We write KCL at these two nodes:
v1 v2 v1 0
+
= is1
R2
R1
v2 0 v2 v1
KCL 2 : is2 + i3 + i2 = 0
+
= is2
R3
R2
KCL 1 : i2 + i1 is1 = 0
The above are two equations in two unknowns (node voltages v1 and
v2 ). Note that a circuit with 5 elements which, in principle, has 10 circuit
variables and 10 equations to be solved is reduced to 2 equations in 2
unknowns. The previous equations can be rearranged as:

1
1
1
+
+ v2
= is1
v1
R1 R2
R2

1
1
1
+ v2
= is2
v1
+
R2
R2 R3
35
or using the denition of conductance, G = 1/R, we have:

(G1 + G2 )v1 G2 v2 = is1 G2 v1 + (G2 + G3 )v2 = is2
In the matrix form, the equations are:

v1
i1s1
(G1 + G2 ) G2
=
G2
G2 + G3
v2
is2
This can be generalized to any circuit with N nodes (let n = N 1).
The node-voltage equation can be written as
Gv = is
where
G11 G12 . . . G1n

G
G22 . . . G2n
21
G = ..
..
..
.. ,
.
.
.
.
Gn1 Gn2 . . . Gnn

v1
v
2
v = .. ,
.
vn
is1
i
s2
is = ..
.
isn
Matrix G is called the conductance matrix. If the KCLs are written as

sum of current exiting each node, the diagonal elements of this matrix will
be positive and the o-diagonal elements will be negative. The diagonal
elements Gii are equal to the sum of all conductances connected to node i.
The matrix is symmetric, i.e. Gij = Gji . These o-diagonal elements, Gij ,
are sum of conductances directly connecting nodes i and j. Vector v is the
array of node voltages. Vector is is the array of current sources. Element
isi is net equivalent of source currents entering the node.
Because of the above features of the node-voltage equations, it is possible
to directly write the node voltage equations in the matrix form of Gv = is
and follow the above rules to construct matrix G and vectors v and is . Node
voltage method is a powerful and simple method.
Example: Consider the circuit in Figure 1.33. Find i.
The circuit has four nodes. Assigning the node at the bottom as the
reference node. Using the rules to construct matrix G and vectors v and
36
5A
v1
7A
12
v2
v3
12
17A
Figure 1.33: Another example for node-voltage method.
is , we obtain
4 1 0
1
G=
1 6 2 ,
12
0 2 7

v1
v = v2 ,
v3

2
is = 5
17
From
Gv = is
we obtain the elements of vector v as v1 = 12V, v2 = 24V, and v3 = 36V.
Using Ohms law, we get
i=
v2 v3
= 2A
6
Circuits with Voltage Sources

Node voltage method as outlined above has to be modied for circuits
with voltage sources as the i-v characteristics of IVS does not specify the
current through it. This modications can be explained in the context of
the example below:
Example: Consider the circuit in Figure 1.34. Find v.
The circuit has four nodes. We will assign the node at the bottom as
the reference node and we have three node voltages as unknowns. If we did
not have voltage sources, we would write 3 KCL at the 3 nodes to arrive at
v1
6K
v2
2K
-
3V
v3
+
20V
37
4K
6mA
Figure 1.34: Circuit with voltage sources.
3 equations in 3 unknowns. For this circuit, writing KCLs is not useful as

the i-v characteristics of IVS does not specify the current through it. We
note that there are two possible conguration for voltage sources:
a) Voltage source is connected to the reference node (such as the 20-V IVS
in the circuit above). If we write the voltage across the IVS in terms of the
node voltages:
v1 0 = 20 v1 = 20V
In this case, we do not need to write any KCL. The value of node voltage
can be derived directly from the i-v characteristics of the IVS.
b) Voltage source is connected between two nodes that are not reference
nodes (such as the 3-V IVS in the above circuit). We need to write two
equations for nodes 2 and 3 (one for each). KCLs for those nodes (two
equations) cannot be used as both contain the current through IVS that is
unknown. One equation is the voltage across the IVS in terms of the node
voltages:
v2 v3 = 3
The second equation can be found by noting that as charge cannot be
accumulated in any part in the circuit, KCL should apply to any portion
of the circuit. Consider a portion of the circuit including nodes 2, 3 and
the 3 V voltage source. We call this portion a supernode (see Figure 1.35).
KCL should be valid for a supernode. This gives
v3
v2
v2 v1
+
+
=0
KCL for supernode 2 & 3 : 6 103 +
4 103 2 103 6 103
72 + 3v3 + 6v2 + 2(v2 v1 ) = 0
2v1 + 8v2 + 3v3 = 72
38
We now have three equations in three unknowns:
= 20
v1
= 3
v2 v3
2v1 + 8v2 + 3v3 = 72

which gives v1 = 20V, v2 = 11V, and v3 = 8V. Therefore problem unknown
v = v3 0 = v3 = 8V.
Supernode
v1 =20V
6K
v2
3V
v3=v2-3
+
2K
20V
4K
6mA
Figure 1.35: Circuit with supernode.
Noting v1 = 20, we can write the value of the node voltage on the circuit.
Also, using v2 v3 = 3, we can arrive at v3 = v2 3 as shown in the gure.
Now, we only have one unknown, v2 , which can be found by writing KCL
on the supernode:
KCL for supernode 2 & 3 : 6 103 +
v3 3
v2
v2 20
+
+
=0
4 103 2 103 6 103
The above is one equations in one unknown which can be solved rapidly.
39
Notes:
1. Voltage sources actually simplify the node voltage equations. By
using the i-v characteristic equations of IVS on the circuit diagram, we can
reduce the number of unknown node voltages (thus, number of equations
to be solved) to N 1 NIV S , where NIV S is the number of IVS.
2. The situation is much simpler for voltage sources that are connected
between a node and ground (such as the 20-V source in the above circuit).
As the location of the reference node is arbitrary, a good rule is: Choose
the reference node as the node with maximum number of voltage sources
attached to it.
3. While we do not need to write KCL at the nodes connected to a
voltage node in order to nd the node voltages, we need to use KCL if we
want to nd the current through a voltage source. For example, in the
circuit above, current i can be found by writing KCL at node 3 (after we
have found all node voltages):
i 6 103 +
v2 3
= 0
4 103
i = 4mA
Summary for Node Voltage Method
1. Identify nodes and supernodes.
a. Choose the reference node as the one with maximum number of
voltage sources attached to it.
b. Use i-v characteristic equations of IVS to nd node voltage values
and reduce the number of unknowns.
2. Write KCL at each node or supernode.
3. Solve node-voltage equations.
4. Calculate problem unknowns from node voltages. If you need to calculate the current in a voltage source you may have to write KCL at a
node connected to that voltage source.
40
1.3.2
Mesh-Current Method
The mesh-current method is analog of the node-voltage method. We solve

for a new set of variables, called mesh currents, that automatically satisfy
KCLs. As such, mesh-current method reduces circuit solution to writing a
bunch of KVLs.
Note: Mesh-current method only works for planar circuits: circuits that
can be drawn on a plane without any elements or connecting wires crossing
each other. Note that in some cases a circuit that looks non-planar can be
made into a planar circuit by moving some of the connecting wires.
i1
i1
+
R2
R1
i2
i3
i1
i2
R3
vs1
vs2
i3
i1
i2
Figure 1.36: Example of mesh-current method.
Mesh-current method is best explained in the context of example circuit

in Figure 1.36.
A mesh is dened as a closed path (a loop) that contains no closed path
within it.
Mesh current is the current that circulates in the mesh i.e.,
a. If an element is located on a single mesh (such as R1 , R2 , vs1 , and vs2 )
it carries the same current as the mesh current,
b. If an element is located on the boundary of two meshes (such as R3 ),
it will carry a current that is the algebraic sum of the the two mesh
currents:
i3 = i1 i2
i3 = i2 i1
41
In this way, KCLs are automatically satised. In addition, as we can

write current in each element in terms of mesh currents, we can use i-v
characteristics of element to write voltage across each element in terms on
mesh currents. Therefore, we need only to write KVLs in terms of mesh
currents. In the circuit above, KVLs give:
Mesh 1 : R1 i1 + R3 (i1 i2 ) vs1 = 0 (R1 + R3 )i1 R3 i2 = vs1
Mesh 2 : R3 (i2 i1 ) + R2 i2 + vs2 = 0 R3 i1 (R2 + R3 )i2 = vs2
or in matrix form,

i1
vs1
R3
R1 + R3
=
R3
R2 + R3
i2
vs2
Ri = vs
which is similar in form to matrix equation found for node-voltage method.

i is the vector of mesh currents (unknowns), vs is the vector of independent
voltage sources, and R is the resistance matrix and is symmetric. The
diagonal element, Rjj , is the sum of resistance around mesh j and the odiagonal elements, Rjk , are the sum of resistance shared by meshes j and
k.
Example: Consider example in Figure 1.37. Find i and v.
+ v i1
i2
i1
6
+
9V
-
16V
-
i2
+
-
6V
Figure 1.37:
method.
Another example of mesh-current
Using mesh-current method:

Mesh 1 : 2i1 + 9 + 3(i1 i2 ) 16 = 0
Mesh 2 : 6i2 + 6 + 3(i2 i1 ) 9 = 0
we get, i1 = 2A, i2 = 1A.
42
The problem unknowns, i and v can now be found from the mesh currents:
i = i1 i2 = 1A
v = 2i1 = 4A
Mesh-Currents Method for Circuits with Current Sources
1A
i1
2
+
10V
-
i2
2
2A
i3
6
Figure 1.38: Mesh-current method for circuit

with current sources.
Because of i-v characteristics of a current source does not specify its

voltage, we have to modify mesh-current method. This is best seen in the
example in Figure 1.38.
From the circuit, we note:
If a current source is located on only one mesh (1-A ICS i1 in the
circuit), the mesh current can be directly found from the current source
and we do not need to write any KVL: i1 = 1A
If a current source is located on the boundary between two meshes (2A ICS in the circuit), KVL on these meshes (mesh 2 or 3 in the above
circuit) contain the voltage across the 2-A ICS which is unknown. We
need two equations to substitute for the two KVLs on meshes 2 and 3
that are not useful now.
The rst equation is found from the i-v characteristics of the current source (its current should be 2 A): i3 i2 = 2A
43
The second equation can be found by noting that KVL can be written over any closed loop. We dene a supermesh as the combination of two meshes which have a current source on their boundary
as shown in the Figure 1.39. While KVL on mesh 2 or on mesh
3 both include the voltage across the 2-A current source that is
unknown, KVL on the supermesh does not include that, we have:
Supermesh 2 and 3 : 2(i2 i1 ) + 2(i3 i1 ) + 6i3 10 = 0
4i1 + 2i2 + 8i3 = 10
1A
i1= -1
2A
+
10V
i2=i3-2
i3
Figure 1.39: Supermesh for circuit with current

sources.
Totally, we have
= 1
i1
= 2
i i2
3
4i1 + 2i2 + 8i3 = 10
which results in i1 = 1A, i2 = 1A, and i3 = 1A.

Summary for Mesh-Current Method
1. Check if circuit is planar.
2. Identify meshes, mesh currents, and supermeshes.
a. Rearrange the circuit if possible to position current source on a
single mesh.
44
b. Use i-v characteristic equations of ICS to nd mesh currents and

reduce the number of unknowns.
3. Write KVL at each mesh and supermesh.
4. Solve for mesh currents.
5. Calculate problem unknowns from mesh currents. If you need to calculate the voltage across a current source you may have to write KVL
around a mesh containing the current source.
6. For consistency and elimination of errors, always mark all mesh currents
in clockwise direction and write down KVLs in the same direction.
Comparison of Node-Voltage and Mesh-Current Methods
Node-voltage and mesh-current are powerful methods that simplify circuit
analysis substantially. They are methods of choice in almost all cases (except for very simple circuits or special circuits). Examination of the circuit
can also tell us which of the two methods are best suited for the circuit
at hand. We always want to reduce the circuit equations into the smallest number of equations in smallest number of unknowns. The number of
equations from node-voltage method, NN V and mesh-current method, NM C
are given by:
NN V = Nnode 1 NV S
NM C = Nmesh NCS
where Nnode and Nmesh are number of nodes and number of meshes, respectively; NV S and NCS are numbers of voltage and current sources, respectively. Thus, always inspect the circuit, nd NV S and NCS , and proceed
with the method that results in the smallest number of equations to solve.
Note: We need to check to ensure that the circuit is a planar circuit. If
it is not one cannot use mesh-current method and should use node-voltage
method.
1.3.3
Principle of Superposition
If a linear circuit is driven by more than one independent source, the response of the circuit can be written as the sum of the responses of the
circuit to individual sources with all other sources having their strength set
to zero. Note that a source is set to zero does not mean removing it.
45
Example: Consider the circuit in Figure 1.40(a). Find v by superposition.

Because we have two independent sources, we rst set the current source
to zero to obtain the circuit in Figure 1.40(b), and then set the voltage
source to zero to arrive at the circuit in Figure 1.40(c). By superposition,
v = va + vb .
Figure 1.40(b) is a voltage divider circuit and va can be written down
directly as
5
15 = 5V
va =
5 + 10
Figure 1.40(c) is a current divider circuit and current i can be written
down directly as
10
3 = 2A
i=
5 + 10
vb = 5i = 10V
Thus,
v = va + vb = 5 10 = 5V
Note: Using superposition results in slightly simpler circuits (one element is replaced with either a short or open circuit) but more circuits.
In general superposition requires more work than node-voltage or meshcurrent methods. Superposition is used:
a. If sources are fundamentally dierent (e.g., DC and AC sources). In this
case superposition may be the only choice.
b. If circuit is repetitive such that circuits resulting from applying superposition look identical and, thus, we need only to solve one circuit.
1.3.4
Reduction of Two-Terminal Subcircuits to Thevenin Form
We used the equivalency of Norton and Thevenin forms in circuit reduction.

Recall our discussion of equivalent elements and subcircuits. We can replace
any two-terminal subcircuit with another one as long as they have the
same i-v characteristics. It is shown that the i-v characteristics of any
two-terminal element containing linear elements is in the Thevenin form of
v = v T RT i
46
10
+
v
+
3A
15V
-
1A
(a)
10
+
v
+
15V
(b)
10
+
v
3A
(c)
Figure 1.40: Example of superposition. (a) Original
circuit. (b) Setting the current source to zero: open
circuit. (c) Setting the voltage source to zero: short
circuit.
47
Consider the two-terminal subcircuit and its Thevenin equivalent (they

have exactly same i-v characteristics) in Figure 1.41. Let the current i = 0
and calculate v, i.e., calculate the voltage across the terminals of the subcircuit, while the terminal are open circuit. This voltage is called the open
circuit voltage, voc . Examination of the Thevenin form shows that if i = 0,
then vT = voc .
RT
i=0
i=0
+
+
+
Subcircuit
vT
v = voc
v = voc
(a)
Subcircuit
i = isc
i = isc
+
iN
v=0
-
RN
v=0
-
(b)
Figure 1.41: Two-terminal subcircuit and its Thevenin/Norton equivalent
circuits. (a) Thevenin equivalent circuit. (b) Norton equivalent circuit.
Next, consider the two-terminal subcircuit and its Norton equivalent

(they have exactly same i-v characteristics). Let the voltage v = 0 and
calculate i, i.e., calculate the current, while the subcircuit terminals are
shorted. This current is called the short circuit current, isc . Examination
of the Norton form shows that if v = 0, iN = isc .
Lastly, if one set all of the sources in the subcircuit to zero, the remaining
circuit should be equivalent to the Thevenins resistance, RT .
48
Finding Equivalent Thevenin/Norton Forms

Compute two of the following quantities by solving the appropriate circuits:
vT = voc , iN = isc , and RT by setting zero the sources. The third parameter
is found from vT = RT iN .
Example: Find Thevenin equivalent of the subcircuit in Figure 1.42.
5
+
25V
-
20
3A
Figure 1.42: Example of Thevenin equivalent.
Thevenins Theorem: Find two of the following three: RT , voc , and isc .
a) Find RT by setting the sources to zero (see Figure 1.43)
5
+
25V
-
20
20
3A
3A
5//20 = 4
Figure 1.43: Circuit for nding RT .
RT = 8
49
b) Find vT = voc (set i = 0)

25V
voc
i=0
v
+
25V
-
20
3A
voc
-
Figure 1.44: Circuit for nding vT = voc .
Using node-voltage method (note that the voltage drop across the 4
resistor is zero).
voc 25 voc 0
+
3 = 0
5
20
vT = voc = 32V
c) Find iN = isc (set v = 0)
25V
v2
4
+
+
25V
-
20
3A
isc
v=0
-
Figure 1.45: Circuit for nding iN = isc .
Using node-voltage method:

v2 0
v2 25 v2 0
+
3+
= 0
5
20
4
v2 = 16V
v2
iN = isc =
= 4A
4
50
1.3.5
Real Sources
Rs
Ideal source
v = vs
R eal
source
vs
vs
is
v
-
(a)
Rs
v
-
(b)
Figure 1.46: Real source. (a) i-v characteristic of real source with a straight line. (b)
Circuit model of a real source.
In an ideal voltage source, the voltage is constant no matter what current

is drawn from the source. In a real, practical voltage source (like a battery),
however, the output voltage typically decreases as more and more current
is drawn, as is shown in the Figure 1.46. Typically, a real source is rated
for currents below a current i which corresponds to a voltage v 95%vs
(region near vs in the gure). For this region, it is a good approximation
to model the i-v characteristics of a real source with a straight line. The
equation of this line is (using active sign convention):
v = v s Rs i
The above approximate i-v characteristics of a real source is a Thevenin
form and, therefore, a real source can be modeled with ideal voltage source,
vs , and a resistance Rs . Rs is called the internal resistance of the source
(note that it is not a real resistor inside the real source)and is typically
small (an ideal voltage source has Rs = 0). The same arguments can be
applied to real current sources. An approximate model for a real current
source is in Norton form. Rs is again the internal resistance of the source
(and again, it is not a real resistor inside the real source). For a real
current source, Rs is typically large (an ideal current source has Rs ).
1.4
Dependent or Controlled Sources
Most analog electronic devices include ampliers. These are four-terminal

devices (two input and two output terminals). The voltage or current in
1.4. DEPENDENT OR CONTROLLED SOURCES
51
the output terminals are proportional to voltage or current of the input

terminals. We need a new circuit element in order to model ampliers.
These elements are controlled or dependent sources. There are four type of
controlled sources as shown in Figure 1.47.
i1
+
+
v1
mv1
ni1
(a )
(b)
i1
qi1
pv1
v1
(c )
(d)
Figure 1.47: Controlled Sources. (a) Voltage-controlled voltage source.

(b) Current-controlled voltage source. (c) Voltage-controlled current
source. (d) Current-controlled current source.
Note that the element located in the input with the controlling current
or voltage can be any element: a short circuit, an open circuit, or a resistor.
Controlled sources behave similar to ideal (or independent) sources. For
example, in the voltage-controlled voltage source in the above gure, the
output voltage is v1 no matter what current is drawn from the circuit. All
analysis method developed so far (KVL and KCL, node-voltage and meshcurrent methods, superposition, etc.) can be used for circuits containing
controlled sources, and by treating the controlled source similar to an ideal
source. In node-voltage and mesh-current methods, however, we need to
write an auxiliary equation which relates the controlling parameter to
node-voltage or mesh-current methods as is seen in the examples below.
52
Rc
Rs
io
ix
Rp
vs
-
+
-
rix
RL
vo
-
Figure 1.48: Example of circuit with a currentcontrolled voltage source.
Example: Find vo in Figure 1.48, using KVL and KCL:

vs
RS + RP
rix
KVL : RC io + RL io + rix = 0 io =
RC + RL
KVL : RS ix + RP ix vs = 0 ix =
Substituting for ix from the rst equation in the second and noting
vo = RL io , we get:

rRL
vo = RL io =
vs
(RS + RP )(RC + RL )
Example: Consider the circuit in Figure 1.49. Find v1 using node-voltage
method.
8i 1
-
16V
-
1 .2 5 A
+
+
v1
0.75v1
i1
Figure 1.49: Example of node-voltage method with controlled sources.
53
Circuit has 5 nodes and 2 voltage sources (one independent and one
controlled voltage source). Thus, the number of node voltage equations is
NN V = 5 1 2 = 2. We choose the reference node to be the one with
most voltage sources attached to it. Then, we get directly
vA = 16
vC = 8i1
8i 1
-
16 V
-
vA = 16V
vC = 8i1
vB
2
+
1. 2 5A
v1
0.75v1
i1
vD
Figure 1.50: Node voltages of the circuit in Figure 1.49.
We then proceed with writing KCL at the other two nodes:

Node vB :
Node vD :
vB 8i1 vB vD
+
1.25 = 0
2
4
3vB vD = 5 + 16i1
vD 0 vD 16 vD vB
+
+
0.75v1 = 0
4
2
4
vB + 4vD = 32 + 3v1
Two above equations are two equations in two unknowns (vB and vD ).
But they also contain the control parameters i1 and v1 . We need to write
two auxiliary equations relating these control parameters to our node
voltages:
i1 =
vD 16
,
2
v1 = vB 8i1
v1 = +vB 4vD + 64
54
We now substitute for control parameters i1 and v1 in our node-voltage

equations to get:
3vB 9vD = 123
3vB vD = 5 + 8(vD 16)
vB + 4vD = 32 + 3(+vB 4vD + 64) 4vB + 16vD = 224
which can be solved to nd the node voltages vB = 4V and vD = 15V. The
control parameters are: i1 = 0.5A and v1 = 8V.
Thevenin Equivalent of Subcircuits with Controlled Sources
Two-terminal subcircuits containing controlled sources can be reduced to
Thevenin form by the method discussed in previous section. The method
was to nd two of three parameters: RT (by setting independent sources
to zero), vT = voc and iN = isc .
Example: Find the Thevenin equivalent of this subcircuit in Figure 1.51.
4i
i
2K
32V
6K
Figure 1.51: Example of Thevenin equivalent of

subcircuit with controlled source.
Finding voc :
Since the circuit is simple, we proceed to solve it with KVL and KCL
(noting that i = 0) :
i1 = 0
KCL : i1 + i + 4i = 0
i2 = 0
KCL : i2 4i + i1 = 0
KVL : 32 + 2 103 i2 + 6 103 i1 + voc = 0
vT = voc = 32V
55
4i
i =0
2K
i2
32V
6K
i1
voc
Figure 1.52: Circuit for nding voc .
Finding isc
Since the circuit is simple, we proceed to solve it with KVL and KCL:
KCL : i1 + i + 4i = 0
i1 = 5isc
KCL : i2 4i + i1 = 0
i2 = isc
3
3
KVL : 32 + 2 10 isc + 6 10 isc = 0
iN = isc = 4mA
4i
i =i s c
2K
+
32V
i2
6K
i1
v=0
Figure 1.53: Circuit for nding isc .
Therefore, vT = 32 V, iN = 4mA, and RT = vT /iN = 8K.

While nding voc and isc is preferred method for most circuits, in some
cases, the Thevenin equivalent of the subcircuit is only a resistor (you will
nd voc = 0 and isc = 0) , or only a voltage source (you will nd voc = 0
but nding isc leads to inconsistent or illegal circuits), or only a current
source (you will nd isc = 0 but nding voc leads to inconsistent or illegal
56
circuits). For these cases, one has to either nd RT directly and/or directly
nd i-v characteristics of the subcircuits as is shown below for the circuit
of previous example.
Finding RT
4i
i
2K
i2
6K
i1
v
-
Figure 1.54: Circuit for nding RT .
To nd RT , we set all independent sources in the circuit to zero. The

resulting circuit cannot be reduced to a simple resistor by series/parallel
formulas. This is why nding voc and isc is the preferred choices for subcircuits containing controlled sources. We can nd RT by attaching an ideal
voltage source with a known voltage of v and calculate i. Since the subcircuit should be reduced to a resistor (RT ), we should get i = v/(constant)
where the constant is RT . (Negative sign comes from active sign convention
used for Thevenin subcircuit).
Since the circuit is simple, we solve it with KVL and KCL:
KCL : i1 + i + 4i = 0
i1 = 5i
i2 = i
KCL : i2 4i + i1 = 0
KVL : 0 + 2 103 i + 6 103 i + v = 0
i=
v
8 103
Therefore, RT = 8 103 = 8K.

Note that we could have attached an ideal current source with strength
of i to the problem, proceeded to calculate v, and obtain v = 8 103 i.
1.5. LINEAR AMPLIFIERS AND OPERATIONAL AMPLIFIERS
57
Finding i-v Characteristics Equation:

As mentioned above, in some cases, we have to directly nd the i-v characteristics equation in order to nd the Thevenin equivalent of a subcircuit.
The procedure is similar to nding RT .
4i
i
2K
+
i2
6K
32V
i1
Figure 1.55: Circuit for nding i-v characteristics

equation.
Attach an ideal voltage source to the circuit. Assume v is known and

proceed to calculate i in terms of v. (Alternatively, one can attach an
ideal current source, assume i is known and nd v in terms of i). The
nal expression should look like v = vT iRT and vT and RT can be read
directly.
Since the circuit is simple, we proceed to solve it with KVL and KCL:
i1 = 5i
KCL : i1 + i + 4i = 0
i2 = i
KCL : i2 4i + i1 = 0
KVL : 32 + 2 103 i + 6 103 i + v = 0
v = 32 8 103 i
which is the characteristics equation for the subcircuit and leads to vT = 32V,
RT = 8 103 = 8K, and iN = vT /RT = 4mA.
1.5
Linear Ampliers and Operational Ampliers
Ampliers are two-port networks in which the output voltage or current is

directly proportional to either input voltage or current. Four dierent kind
of ampliers exits:
Voltage amplier: Av = Vo /Vi : voltage gain (constant)
58
Current amplier: Ai = Io /Ii : current gain (constant)

Transconductance amplier: Gm = Io /Vi : constant
Transresistance amplier: Rm = Vo /Ii : constant
Figure 1.56: Voltage amplier model.
In this section, we will focus on voltage ampliers. Voltage ampliers

can be accurately modelled with three circuit elements as in Figure 1.56.
Ri and Ro are, respectively, input and output resistances. A good voltage
amplier has a large input resistance, and a small output resistance. An
ideal voltage amplier has Ri and Ro 0.
1.5.1
Practical Ampliers
Amplier Saturation
Ampliers do not create power. Rather, they act as a valve adjusting
the power ow from the power supply into the load according to the input
signal. As such, the output voltage amplier cannot exceed the power
supply voltage (it is usually lower because of voltage drop across some active
elements). The fact that the output voltage of a practical amplier cannot
exceed certain threshold value is called saturation. A voltage amplier
behaves linearly, (i.e., Vo /Vi = Av : constant), as long as the output voltage
remains below the saturation voltage, Vsat
Vsat < Vo < Vsat
Note that the saturation voltage, in general, is not symmetric, i.e.,
Vsat,1 < Vo < Vsat,2 . For an amplier with a given gain, Av , the above
range of Vo translate into a certain range for Vi
Vsat < Vo < Vsat
59
since Vo = Av Vi , then
Vsat
Vsat
< Vi <
Av
Av
i.e., any amplier will enter its saturation region if Vi is raised above certain
limit.
Figure 1.57: Amplier output clipped when amplier

is not in linear region.
Amplier Bandwidth
Typically, ampliers work in a certain range of frequencies. Their gain,
Av = Vo /Vi drops outside this range. The frequency response of amplier
is plotted as in Figure 1.58, where fL , fH are 3-dB cut-o frequencies. The
bandwidth of the amplier is dened as BW = fH fL .
In terms of frequency response, voltage ampliers are divided into two
categories. (1) AC ampliers which only amplify AC signals. (2) DC ampliers which amplify both DC signals and AC signals up to a certain
frequency.
60
|A |
Am
Mid-frequency range
A m / 1 .4
High-frequency
range
Low-frequency range
fL
f
fH
Figure 1.58: Frequency response of amplier.
Rise Time
In an ideal amplier, if the input voltage is a unit step function, the output
voltage will also be a unit step function as shown. A practical amplier
cannot change its output instantaneously if the input changes suddenly. It
takes some time (a short but nite time) for the amplier output voltage to
reach its nominal level. The maximum rate of change in the output voltage
is called the rise time.
Figure 1.59: Rise time of practical

amplier.
1.5.2
Feedback
Not only a good amplier should have sucient gain, its performance should
be insensitive to environmental and manufacturing conditions, should have
61
a large Ri , a small Ro , a suciently large bandwidth, etc. It is easy to make

an amplier with a very large gain. A typical transistor circuit can easily
have a gain of 100 or more. A three-stage transistor amplier can easily get
gains of 106 . Other characteristics of a good amplier is hard to achieve,
such as, the gain of amplier changes with operating temperature. The
system can be made to be insensitive to environmental and manufacturing
conditions by the use of feedback.
Principle of Feedback
The input to the circuit is modied by feeding a signal proportional to
the output value back to the input. There are two types of feedback:
Negative feedback: As the output is increased, the input signal is
decreased and vice versa. Negative feedback stabilizes the output to
the desired level. Linear system employs negative feedback.
Positive feedback: As the output is increased, the input signal is increased and vice versa. Positive feedback leads to instability.
We will explore the concept of feedback in the context of operational
ampliers.
1.5.3
Operational Ampliers as Linear Ampliers
Figure 1.60: OpAmps as a dierence amplier.
Operational ampliers (OpAmps) are general purpose voltage ampliers

employed in a variety of circuits. OpAmps are DC ampliers with a
62
very large gain, (105 to 106 ), high input impedance (> 1 10M), and
low output resistance (< 100). They are constructed as a dierence
amplier, i.e., the output signal is proportional to the dierence between
the two input signals (see Figure 1.60)
Vo = A0 Vd = A0 (Vp Vn )
In Figure 1.60, Vs and Vs are power supply attachments. They set the
saturation voltages for the OpAmp circuit. Power supply ground should
also be connected to the OpAmp ground. + and - terminals of the OpAmp
are called, respectively, non-inverting and inverting terminals.
OpAmp Models
Figure 1.61: OpAmp models.
Because Ri is very large and Ro is very small, ideal model of the OpAmp
assumes Ri and Ro 0. Ideal model is usually a very good model
for OpAmp circuits. Very large input resistance also means that the input
current into an OpAmp is very small:
Ip In 0
An important feature of OpAmp is that because the gain is very high,
the OpAmp will be in the saturation region without negative feedback. For
example, take an OpAmp with a gain of 105 and Vsat = 15V. Then, for
OpAmp to be in linear region, Vi 15 105 = 150V (a very small
value). OpAmps are never used by themselves. They are always part of a
circuit which employ either negative feedback (e.g., linear ampliers, active
lters) or positive feedback (e.g., comparators). Examples below shows
several OpAmp circuits with negative feedback.
1.5.4
63
Inverting Amplier
Figure 1.62: Inverting amplier.
The rst step in solving OpAmp circuits is to replace the OpAmp with
its circuit model (ideal model is usually used).
Vp = 0,
Vo = A0 Vd = A0 (Vp Vn ) = A0 Vn
Using node-voltage method and noting In 0:

Vn Vi Vn Vo
+
=0
R1
R2
Substituting for Vn = Vo /A0 and multiplying the equation by R2 , we
have:
R2 Vi
Vo
R2 Vo
Vo = 0
A0 R1
R1
Ao

Vo
R2
1
=
Vi
R1
1 + 1/A0 + R2 /(A0 R1 )
Since OpAmp gain is very large, 1/A0 1. Also if R2 and R1 are chosen
such that their ratio is not very large, R2 /R1 A0 or R2 /(A0 R1 ) 1,
then the voltage gain (voltage transfer function) of the OpAmp is
Vo
R2
=
Vi
R1
The circuit is called an inverting amplier because the voltage gain is
negative. The negative sign means that there is 180 phase shift between
input and output signals.
Note that the voltage gain of the inverting amplier is independent
of the OpAmp gain, A0 , and is only set by the values of the resistors R1
64
and R2 . While A0 is quite sensitive to environmental and manufacturing

conditions, the resistor values are quite insensitive and, thus, the gain of
the system is quite stable. This stability is achieved by negative feedback,
the output of the OpAmp is connected via R2 to the inverting terminal
of OpAmp. If Vo increases, this resistor forces Vn to increase, reducing
Vd = Vp Vn and Vo = A0 Vd , and stabilizes the OpAmp output.
An important rule of OpAmp circuits with negative feedback is that
because the OpAmp is NOT saturated, Vd = Vo /A0 is very small (because
A0 is very large). As a result, Vd 0 Vn Vp
Note: For OpAmps circuits with negative feedback, the OpAmp adjusts
its output voltage such that Vd 0 or Vn Vp . This rule is derived by
assuming A . Thus, Vo cannot be found from Vo = A0 Vd = 0 =
indenite value.
Input and Output Resistances of Inverting Amplier
From the circuit,
Vi 0
Vi
Ri =
= R1
R1
Ii
The input impedance of the inverting amplier circuit is R1 (although
input impedance of OpAmp is innite). The output impedance of the
circuit is zero because Vo is independent of RL (Vo does not change when
RL is changed).
Ii =
1.5.5
Non-inverting Amplier
Vp = Vi
Negative Feedback Vn Vp = Vi
Vn 0 Vn Vo
+
=0
R1
R2
Substituting for Vn = Vi , we get
R2
Vi + Vi Vo = 0
R1
Finally,
Av =
Vo
R2
=1+
Vi
R1
Vi
Ii
Vi
Vo
Ip
Vp
Vo
Vd
AVd
R2
65
Vn
R1
R2
R1
Figure 1.63: Non-inverting amplier.
Input Resistance
Since Ii = Ip = 0. Therefore, Ri .
Output Resistance
Vo is independent of RL , so Ro = 0. Note that Ri and Ro = 0 should be
taken in the context that we are using an ideal OpAmp model. In reality,
the above circuit will have input and output resistances equal to that of
the OpAmp itself.
1.5.6
Voltage Follower
Vi
+
-
Figure 1.64: Voltage follower.
Vo
66
In some cases, we have two-terminal networks which do not match well,

i.e, the input impedance of the later stage is not very large, or the output
impedance of preceding stage is not low enough. A buer circuit is
usually used in between these two circuits to solve the matching problem.
These buer circuit typically have a gain of 1 but have a very large input
impedance and a very small output impedance. Because their gains are 1,
they are also called voltage followers.
The non-inverting amplier above has Ri and Ro = 0 and, therefore, can be turned into a voltage follower (buer) by adjusting R1 and R2
such that the gain is 1
R2
Vo
=1+
= 1 R2 = 0
Vi
R1
So by setting R2 = 0, we have Vo = Vi or a gain of unity. We note that
this expression is valid for any value of R1 . As we want to minimize the
number of components in a circuit, we can remove R1 from the circuit.
1.5.7
Inverting Summer
Rf
V2
V1
R2
R1
Vn
Vp
Vo
+
+
-
AVd
Figure 1.65: Inverting summer.
Vp = 0
Negative Feedback Vn Vp = 0
Vn V1 Vn V2 Vn Vo
+
+
=0
R1
R2
Rf
67
Finally,
Rf
Rf
V1
V2
R1
R2
So, this circuit adds (sums) two signals. An example of the use of this
circuit is to add a DC oset to a sinusoidal signal.
Vo =
1.5.8
Non-Inverting Summer
V2
V1
R2
Vp
R1
Vn
Vo
+
-
AVd
Rf
Rs
Figure 1.66: Non-inverting summer.
Negative Feedback: Vn Vp
Vp V1 Vp V2
+
= 0 Vp
R1
R2
Vn 0 Vn Vo
+
=0
R1
Rf
1
1
+
R1 R2
Vo =
Rf
1+
Rs
V2
V1
+
R1 R2

Vn
Substituting for Vn in the second equation from the rst (noting Vp = Vn ),

we obtain

V2
V1
1 + Rf /Rs
+
Vo =
1/R1 + 1/R2
R1 R2
So, this circuit also adds (sums) two signals. It does not, however, inverts
the signals.
68
1.5.9
Dierence Amplier
Rf
V1
R1
Vn
Vp
V2
R2
Vo
+
+
-
AVd
R3
Figure 1.67: Dierence amplier.
Negative Feedback: Vn Vp
Vp V2 Vp 0
+
=0
R2
R3
Vn Vp =
R3
V2
R2 + R3
Vn V1 Vn Vo
+
=0
R1
Rf
If one choose the resistors such that R3 /R2 = Rf /R1 , then
Vo =
1.5.10
Rf
(V2 V1 )
R1
Current Source
Negative Feedback: Vn Vp = Vs
Vs
Vn
= : constant
R
R
For a xed value of Vs , the current IL is independent of value of RL and
output voltage Vo . As such, this circuit is an independent current source.
The value of the current can be adjusted by changing Vs (for a xed RL .
Therefore, this circuit is also a voltage to current converter or a transconductance amplier.
iL =
1.6. DYNAMIC OR TIME-DEPENDENT CIRCUITS
Vp
Vs
Vn
69
+
IL
+
Vo
RL
-
Figure 1.68: Current source.
Grounded Current Source

Rf
Vs
R1
Vn
Vp
R2
R3
RL
IL
+
Vo
-
Figure 1.69: Grounded current source.
The problem with the above current source is that the load is not
grounded. This may not be desirable in some cases. This circuit here
is also a current source with a grounded load if Rf /R1 = R3 /R2 .
1.6
Dynamic or Time-Dependent Circuits
In this section we discuss circuits that include capacitors and inductors.

The i-v characteristics of these two elements include derivatives or integral
70
of either i or v resulting in time-dependent circuits.

Capacitor
A typical capacitor is made of two parallel conducting plates separated by
an insulator. If we connect a capacitor to a voltage source, electrons travel
from the voltage source to the capacitor and charge one of the plates with
excess electrons. The electric eld from these electrons repel electrons from
the opposing plate and equal but positive charge develops on that plate.
The repelled electrons from the positive plate travel to the source. As such,
the total charge stored in a capacitor is zero but consists of two separated
amount of equal charges with dierent polarity. In this manner, a capacitor
can store energy in the form of electric eld between the two plates, which
is proportional to the voltage between the two plates. For most parallel
plate capacitors, the voltage across the capacitor is directly proportional to
the charge accumulated in the capacitor. The constant of proportionality is
the capacitance. Note that the total charge stored in a capacitor is zero, the
charge mentioned here is the charge on one of the plates. The capacitance,
C is
q
Unit: Farad (F)
C=
v
Most commercial capacitors are in F (106 F), nF (109 F), and pF
(1012 F) range.
As the number of electrons owing into the negative plate is equal to the
number of electrons owing out of the positive plate, the current owing
into one plate of capacitor is exactly equal to the current owing out of the
other plate. The i-v characteristics equation for a capacitor (using passive
sign convention) is
q = Cv
i=
dv
dq
=C
dt
dt
Note that the current i ows as long as the voltage (and charge, q = Cv)
changes in time. At the DC steady-state condition, dv/dt = 0 and capacitor
current, i = 0, and capacitor acts as an open circuit.
Energy stored in a capacitor, W , can be found from the denition of
electric power:
1
W (t) = Cv 2 (t)
2
71
Time dependent circuits either have a time-dependent source (we will

examine circuit with sinusoidal sources later) and/or have a source which
is switched on or o at some time. In all switching circuits, we assume
that switch is thrown instantaneously. This means that for these switched
circuit, the voltage or current waveform will have a discontinuity in time.
For example, consider the circuit shown in Figure 1.70. The switch is
initially closed and is opened at time t0 . For t < t0 , the switch is closed.
By KCL, i = is and a voltage of v = Ris appears across the resistor. For
t > t0 , i = 0 and v = 0. There is a discontinuity in voltage and current at
time t = t0 when the switch is thrown.
t = t0
t < t0
i
Is
+
v
i
R
Is
+
v
-
t > t0
i
Is
+
v
-
Figure 1.70: Switched resistive circuit.
Values of current and voltage at the exact switching time is undened.

The values are known just prior to the throwing the switch and just after
that. We denote that time just prior to throwing the switch as t

0 (t0 = t0
with being innitesimally small). We denote the that time just after
+
throwing the switch as t+
0 (t0 = t0 + ). In this sense, current and voltage
+
are well dened at t
0 and t0 . In the circuit shown, both current and voltage
+
+
have a discontinuity at t = t0 as i(t
0 ) = i(t0 ) and v(t0 ) = v(t0 ).
This behavior is in general true for all resistive circuit. There is a discontinuity in both voltage and current waveforms at the time that the switch
72
is thrown. This is not true for circuits with capacitors (or inductors). For
capacitors, the i-v characteristics equation i = Cdv/dt implies that if there
were a discontinuity in the voltage, an innite amount of current should
ow through the capacitor. Therefore, the voltage waveform across a capacitor should be continuous. This means that if a switch is thrown at time
+
t0 , we should have v(t
0 ) = v(t0 ). Note the capacitor current waveform can
have a discontinuity.
Inductor
A typical inductor is made of a a wire wrapped around a core of magnetic
material. As current ows through the wire, a magnetic eld is produced.
As such, an inductor can store energy in the form of magnetic eld. The
i-v characteristics equation for an ideal inductor is given by Faradays Law
(wire in an ideal inductor has no resistance):
di
Inductance, L. Unit: Henry (H)
dt
using passive sign convention. Energy stored in an inductor, W , can be
found as
1
W (t) = Li2 (t)
2
For inductors, the i v characteristics equation v = Ldi/dt implies that
if there were a discontinuity in the current waveform, an innite amount
of voltage should appear across the inductor. Therefore, the current waveform in an inductor should be continuous. This means that if a switch is
+
thrown at time t0 , we should have i(t
0 ) = i(t0 ). Note the inductor voltage
waveform can have a discontinuity.
v=L
Comments
Circuits containing one or more capacitors and/or inductors are called
dynamic circuits. Voltage and currents in dynamic circuits are, in
general, functions of time, i.e., their value change as time progresses.
All analysis method (KVL, KCL, circuit reduction, node-voltage and
mesh current methods) can be applied to dynamic circuits with no
modication.
Application of the analysis methods to a dynamic circuit results in
a set of linear dierential equations (as opposed to a set of algebraic
73
equations for resistive circuits). Solution of these equations requires

initial conditionsone for each capacitor or inductor.
Two-terminal subcircuits containing capacitors and/or inductors cannot be reduced to Thevenin or Norton forms.
Switched Circuits
Time dependent circuits either have a time-dependent source and/or have a
source which is switched on or o at some time. Most circuits have both.
The voltage source in the circuit is a time-dependent voltage (sinusoidal
voltage with a frequency of 60 Hz). When the switch is o, the voltages
and currents in the is zero. After the switch is thrown, the voltages and
currents rise from zero and after some time start to follow the source voltage
waveform in a steady manner.
The evolution of circuit voltages and current in time, thus, can be divided into two distinct time scales. 1) A transient response starting with
the throw of the switch and 2) A steady response to the time dependent
sources. We will see that the transient response of the circuit is completely
set by the circuit elements themselves and not by the source. This transient response also stops after some prescribed time. Because of these
two distinct responses of a time-dependent circuit, transient and steady
responses can be studied separately.
As an example, consider a switched circuit as in Figure 1.71. The switch
is open at t = t0 . Is is DC current source. It is usually assumed that the
circuit has been in the rst state (t < t0 ) for a long time and the response
of the circuit for t > t0 is needed.
The solution to the circuit for t > t0 leads to a set of linear dierential
equations. These equations require a set of initial conditions (one for each
capacitor and inductor). These initial conditions are found by solving the
circuit for t < t0 and using the no-discontinuity conditions for voltages
across capacitors and currents in inductors at the switching time. So, the
procedure for solving switched circuits with DC sources is:
Solve DC steady-state case for t < t0 . Determine vc (t

0 ) and iL (t0 ).
Use no-discontinuity conditions to nd the initial condition for the
second circuit: vc (t+

0 ) = vc (t0 ) and iL (t0 ) = iL (t0 ).
74
t = t0
Is
t < t0
Is
t > t0
Is
Figure 1.71: Switched circuit with capacitor.
Solve the time dependent circuit for t > t0 .

DC Steady-State Analysis
In a dynamic circuit with DC sources, voltages and currents reach constant
values after long time. This state of the circuit is called DC Steady
State. By denition, all time derivatives vanish at DC steady state condition:
dv
(for any v). Thus, a capacitor acts as an open
For a capacitor, i = C
dt
circuit DC steady state.
di
For an inductor, v = L (for any i). Thus, an inductor acts as a short
dt
circuit in DC steady state.
Therefore, to analyze a circuit in DC steady state condition, we replace
capacitors with open circuit and inductor with short circuit and proceed to
solve the resulting resistive circuit.
Example
This circuit has been in DC steady state and t = t0 the switch is opened at
time t = t0 . Find the initial conditions for the dynamic circuit for t > t0 .
First, we mark vC and iL on the original circuit (those are our initial
t = t0
1A
200
200
iL
+
vc
-
vc
1A
75
50
iL
+
vc
-
50
Figure 1.72: Example of switched circuit with capacitor

and inductor.
76
conditions). We draw the circuit for t0 < t (closed switch) and replace the
capacitor with an open circuit and the inductor with a short circuit. We
then proceed to solve the resistive DC steady state circuit.
The resulting circuit is current divider and iL and vC can be found readily
from current-divider formulas. Alternatively, using node-voltage method,
we have
vc 0 vc 0
+
1 = 0 vc = 40V
50
200
vc 0
= 0.8A
iL =
50
Note that since the circuit is in DC steady state, values of vC and iL are
constant in time. Using the no-discontinuity conditions for vC and iL , we
can nd the initial conditions for the time-dependent circuit as
vc (t = t+
0 ) = 40V
iL (t = t+
0 ) = 0.8A.
Chapter 2
Introduction to Signals and Systems
2.1
Classication of Signals
A signal is a set of information or data. In this lecture, we deal almost

with signals that are function of time, f (t).
A convenient measure of size of signal is
signal energy Ef (if signal is nite: signal amplitude 0, as |t| )
dened as

f 2 (t)dt
Ef =
for a complex valued signal f (t), we have

|f (t)|2 dt
Ef =
or
signal power (if signal energy is innite: signal amplitude does not
0, as |t| ) which is time average of energy (if it exists). We
dene power Pf of f (t) as

1 T /2 2
f (t)dt
Pf = lim
T T T /2
for a complex valued signal f (t), it is

1 T /2
Pf = lim
|f (t)|2 dt.
T T T /2
For periodic signals, the time averaging need to be performed only over
one period. Square root of signal power is rms (root mean square) value
of the signal.
77
CHAPTER 2. INTRODUCTION TO SIGNALS AND SYSTEMS
78
Signals can be classied as follows:

A continuous-time signal is specied for a continuum of time values.
A discrete-time signal is specied for a nite set of time instants.
A analog signal is a signal whose amplitude can take on any value over
a continuum. On the other hand, a signal whose amplitude can take
on only a nite number of values is a digital signal.
A periodic signal f (t) is dened by the fact that f (t) = f (t + T0 ),
for some T0 . A signal is aperiodic if it is not periodic. An everlasting
signal starts at t = and continues to t = . A causal signal is a
signal that is zero for t < 0. Therefore, periodic signals are everlasting
signals.
A signal with nite energy is an energy signal. A signal with nite and
nonzero power (mean square value) is a power signal.
A signal whose physical description is known completely in a mathematical or graphical form is a deterministic signal. A random signal or
stochastic signal is known only in terms of its probabilistic description
such as mean value, mean square value, etc., rather than its mathematical or graphical form. In this lecture we deal with deterministic
signals.
2.2
2.2.1
Useful Signal Operations

Time Shifting
A signal f (t) delayed by T seconds (right-shifted) is given by f (t T ). A

signal f (t) advanced by T seconds (left-shifted) is given by f (t + T ).
2.2.2
Time Scaling
The compression or expansion of signal in time is known as time scaling.

A signal f (t) time-compressed by a factor a(a > 1) is given by f (at). A
signal f (t) time-expanded by a factor a is given by f ( at ). The same signal
time-inverted is given by f (t).
2.2.3
Combined Operations
(see p. 66, textbook)
2.3. USEFUL SIGNAL MODELS
2.3
79
Useful Signal Models
The below signal models serve as a basis for representing other signals and
their use can simplify many aspects of the signal and systems.
2.3.1
Unit Step Function
The unit step function u(t) is very useful in representing causal signals
and signals with dierent mathematical description over dierent intervals.
This function is dened by

1 t0
u(t) =
0 t<0
2.3.2
Unit Impulse Function
The unit impulse function (Dirac function) (t) is dened as

(t) = 0
(t)dt = 1
t = 0
The unit impulse can be approximated by exponential pulse, triangular

pulse, or Gaussian pulse.
Multiplication of a Function by an Impulse
Multiplying unit impulse (t) by a function (t) that is continuous at t = 0
(t)(t)
= (0)(t)
(t)(t T ) = (T )(t T )
Sampling Property of Unit Impulse Function
The below property is known as the sampling or sifting property of unit
impulse

(t)(t)
=
(0)
(t) = (0)

(t)(t T ) = (T ) (t T ) = (T )
which states that the area under the product of a function with a unit
impulse is equal to the value of that function at the instant where the
impulse is located (assuming the function to be continuous at the impulse
location).
80
2.3.3
Exponential Function
One of the most important functions in signals and systems is the exponential function est , where s = + j is thecomplex frequency. Therefore
est = et+jt = et (cos t + j sin t)
If s is the conjugate of s, then
es t = etjt = et (cos t j sin t)

Then
et cos t = (est + es t )
2
The exponential function encompasses a large class of signals as
A constant k = ke0t (s = 0)
A monotonic exponential et
A sinusoid cos t ( = 0, s = j)
An exponentially varying sinusoid et cos t (s = j)
2.4
Even and Odd Functions
A function fe (t) is an even function of t if

fe (t) = fe (t)
that is symmetrical about the vertical axis (t = 0).
A function fo (t) is an odd function of t if
fo (t) = fo (t)
that is anti-symmetrical about the vertical axis.
Properties
Even and odd functions have the following properties:
even function odd function = odd function
odd function odd function = even function
2.5. CLASSIFICATION OF SYSTEMS
81
even function even function = even function

Area: The area under an odd function from t = a to a is always zero
regardless of the value of a
a
fo (t)dt = 0.
a
On the other hand, the area under an even function from t = a to a

is two times the area under the same function from t = 0 to a
a
a
fe (t)dt = 2
fe (t)dt.
a
Every signal f (t) can be expressed as a sum of even and odd components because
1
1
f (t) = [f (t) + f (t)] + [f (t) f (t)] = fe (t) + fo (t)
2
2
2.5
Classication of Systems
A system processes input signals to modify them or extract additional information from them to produce output signals (response). A system may
be made up of physical components (hardware) or may be an algorithm
that computes an output signal from an input signal (software).
A system is characterized by its input, its output, and the rules of operation describing its behavior. Example, in electrical engineering, the rules
are the i-v characteristic of resistors, capacitors, inductors, transistors, etc.,
as well as the rules for interconnection such as Kirchhos laws. Using these
rules we derive mathematical equations relating the outputs to the inputs.
These mathematical equations represent a mathematical model of the system.
The study of systems consists of three major areas:
Mathematical modeling
Analysis problem: To determine the system outputs for the given inputs and a given mathematical model of system.
Design problem (or synthesis problem): To construct a system which
produce a desired outputs for the given inputs.
82
(See example on p. 78)

A system processes input signals to produce output signals. The input
is the cause and the output is its eect. In general, the output is aected
by two causes: the internal conditions of the system (such as the initial
conditions) and the external input.
Systems can be classied in several ways:
2.5.1
Linear and Nonlinear Systems
Linear systems are characterized by the linear property, which implies superposition. If several causes (such as various inputs and initial conditions)
are acting on a linear system, the total eect (response) is the sum of the
responses from each cause, assuming that all causes are absent. A system
is nonlinear if it is not linear.
Response of a Linear System
A linear system output for t 0 must be the sum of two components:
rst, the zero-input response component that results only from the
initial condition at t = 0 with the input f (t) = 0 for t 0
and then the zero-state response component that results only from
input f (t) for t 0 when the initial conditions (at t = 0) are assumed
to be zero.
The property of linear systems which permits the separation of an output
into components resulting from the initial conditions and from the input is
called the decomposition property.
(See example 1.9, p. 81)
In general, a system described by a dierential equation of the form
dn1 y
dm f
dm1 f
df
dn y
+
a
+
+
a
y
=
b
+
b
+
+
b
+ b0 f
n1
0
m
m1
1
dtn
dtn1
dtm
dtm1
dt
is a linear system. The coecients ai and bi can be constant or functions
of time.
2.5. CLASSIFICATION OF SYSTEMS
2.5.2
83
Time-Invariant and Time-Varying Systems
The time-invariant systems are characterized by the fact that the system
parameters do not change with time. For example if the coecients ai and
bi of the linear system are constant, then the linear system is linear timeinvariant (LTI) system. If the system parameters are function of time, then
the system is time-varying system.
2.5.3
Memoryless and Memory Systems
For memoryless (instantaneous) systems, the system response at any instant t depends only on the present value (at t) of the input. For systems
with memory (dynamic systems), the system response at any instant t depends not only on the present value (at t) of the input but also on the past
values of the inputs (value before t).
2.5.4
Analog and Digital Systems
A system whose input and output signals are analog is an analog system; a
system whose input and output signals are digital is a digital system.
2.5.5
Continuous-Time and Discrete-Time Systems
Systems whose inputs and outputs are continuous-time signals are continuoustime systems. On the other hand, systems whose inputs and outputs are
discrete-time signals are discrete-time systems.
2.5.6
Lumped-Parameter and Distributed-Parameter Systems
If the dimensions of system elements are small compared to the wavelength

of the signals, we assume that each element is lumped at a single point in
space, and the system is considered as a lumped-parameter system. The
signals under this assumption are function of time only. For such systems,
the system equations require only one independent variable (time). If the
above assumption does not hold (at very high frequency, in systems such
as transmission lines, antennas, waveguides, etc.), the signals are functions
of space and time, such a system is a distributed-parameter system.
84
2.5.7
Causal and Noncausal Systems
A causal (or physical or non-anticipative) system is one for which the output
at any instant t0 depends only on the value of the input f (t) for t t0 .
In other words, the value of the output at present instant depends only
on the past and present values of the input f (t), not on its future values.
A system that violates the conditions of causality is called a noncausal or
anticipative system.
Any practical system operating in real-time (In real-time operations, the
response to an input is simultaneous with the input itself) must be causal.
Noncausal systems are unrealizable in real time (but realizable with some
time delay in response). Noncausal systems with independent variables
other than time (e.g. space) are realizable.
2.6
System Model: Input-Output Description
To construct a system model, we must study the relationships between different variables in the system. In electrical engineering, we must determine
a satisfactory model for i-v relationship of each element (such as Ohms
law) and the various constraints on voltages and currents when several elements are interconnected (such as KCL, KVL). For all these equations,
we eliminate unwanted variables to obtain equations relating the desired
output variable to the inputs.
It is convenient to use a compact notation D for the dierential operator
d
dt
dy
dt
Dy(t)
d2 y
D2 y(t)
2
dt
..
.
For example,
d2 y
df
dy
+
2y(t)
=
,
+
3
dt2
dt
dt
can be expressed as
(D2 + 3D + 2)y(t) = Df (t)
2.6. SYSTEM MODEL: INPUT-OUTPUT DESCRIPTION
We can use the operator

1
to represent integration
D
t
y( )d
1
y(t)
D
85
Chapter 3
Time-Domain Analysis of
Continuous-Time Systems
In this chapter, we consider time-domain method to analyze linear, time
invariant, continuous-time (LTIC) systems.
We consider linear dierential systems. This is the class of LTIC, for
which the input f (t) and the output y(t) are related by linear dierential
equations as
dn1 y
dm f
dm1 f
df
dn y
+b0 f (t) (3.1)
+a
+
+a
y
=
b
+b
+
+b
n1
0
m
m1
1
dtn
dtn1
dtm
dtm1
dt
where all the coecients ai and bi are constant. Using notation D to
represent d/dt, we obtain
(Dn +an1 Dn1 + +a1 D+a0 )y(t) = (bm Dm +bm1 Dm1 + +b1 D+b0 )f (t)
(3.2)
or
Q(D)y(t) = P (D)f (t)
(3.3)
where the polynomials Q(D) and P (D) are
Q(D) = Dn + n1 Dn1 + + a1 D + a0
P (D) = bm Dm + bm1 Dm1 + + b1 D + b0
(3.4)
A system described by equation (3.1) is linear. Therefore, its response

can be expressed as the sum of two components:
Total response = zero-input response + zero-state response
The zero input component is the system response when the input f (t) = 0
86
3.1. ZERO-INPUT RESPONSE
87
so that it is the result of internal system conditions (such as energy storages,

initial conditions) alone. It is independent of the external input f (t).
The zero input component is the system response to the external input
f (t), when the system is in zero state,; that is, all initial conditions and all
initial energy storages are zero.
3.1
Zero-Input Response
The zero-input response y0 (t) is the solution of equation (3.1) when the
input f (t) = 0 so that
n
(D + n1 D
n1
Q(D) = 0
+ + a1 D + a0 )y0 (t) = 0
(3.5)
The general solution of equation (3.5) is given by

y0 (t) = c1 e1 t + c2 e2 t + + cn en t
(3.6)
where c1 , c2 , , cn are arbitrary constant determined by n constraints (the

auxiliary conditions) on the solution; 1 , 2 , , n are characteristic roots
which are the roots of the characteristic equation of the system:
Q() = 0
(3.7)
The polynomial Q() is called the characteristic polynomial of the system

(note that Q() is identical to Q(D) in equation (3.5) with replacing D).
The exponentials
ei t
(i = 1, 2, n)
in the zero-input response are the characteristic modes of the system. There
is a characteristic mode for each characteristic root of the system, and the
zero-input response is a linear combination of the characteristic modes of
the system
3.1.1
Repeated Roots
The solution of equation (3.5) is given in equation (3.6) assumes that n

characteristic roots 1 , 2 , , n are distinct. If there are repeated roots,
88
CHAPTER 3. TIME-DOMAIN ANALYSIS OF CONTINUOUS-TIME SYSTEMS
then from equation (3.5) we obtain

(D )r y0 (t) = 0
(3.8)
the characteristic modes are

et , tet , t2 et , , tr1 et
and the solution is
y0 (t) = (c1 + c2 t + + cr tr1 )et
(3.9)
Generally, for a system with the characteristic polynomial as

Q() = ( 1 )r ( r+1 ) ( n )
(3.10)
the characteristic modes are

e1 t , te1 t , , tr1 e1 t , er+1 t , , en t
and the solution is
y0 (t) = (c1 + c2 t + + cr tr1 )e1 t + cr+1 er+1 t , , cn en t
3.1.2
(3.11)
Complex Roots
For a real system, complex roots must occur in pairs of conjugates if the
coecients of the characteristic polynomial Q() are to be real. The zeroinput response corresponding to this pair of complex conjugate roots is
y0 (t) = c1 e(+j)t + c2 e(j)t
(3.12)
For a real system, the response y0 (t) must also be real. This is possible
only if c1 and c2 are conjugates. Let
c
c
c1 = ej and c2 = ej
2
2
Then
c j (+j)t c j (j)t
+ e e
e e
2
2

c t j(t+)
j(t+)
= e e
+e
2
y0 (t) =
= cet cos(t + )
(3.13)
3.2. ZERO-STATE RESPONSE
89
Therefore, the zero-input response corresponding to complex conjugate

roots j can be expressed in a complex form (3.12) or a real form (3.13).
(See example 2.1, p. 108 and example 2.2, p. 111)
3.2
Zero-State Response
In this section we determine the zero-state response of LTIC response. This

is the system response y(t) to an input f (t) when the system is in zero-state
(when all initial conditions are zero).
Linear Systems Response to Arbitrary Input
We use the superposition principle and the impulse (Dirac) function to
derive a linear systems response to arbitrary input f (t). The process is
illustrated in Figure 3.1.
(t)
h(t)
(t nt)
h(t nt)
(3.14)
f (n )(t n ) [f (n ) ]h(t n )

input
output
The last pair in (3.14) represents the system response to only one of
the impulse components of f (t). The total response y(t) is obtained by
summing all such components as depicted in Figure 3.1(d).
lim

n=
f (n )(t n )

The input f (t)
lim
f (n )h(t n )
n=

The output y(t)

(3.15)
Both sides in (3.15), by denition, are integrals given by

f ( )(t )d
f ( )h(t )d

f (t)
y(t)
(3.16)
90
Figure 3.1: Finding the system response to an arbitrary input f (t).
91
Therefore, the zero-state response y(t) to the input f (t) is given by

y(t) =
f ( )h(t )d
(3.17)
Note that the assumption used to derive (3.17) is linear and time-invariant
system.
3.2.1
Convolution Integral
The convolution integral of two functions f1 (t) and f2 (t) denoted by f1 (t) f2 (t)
is dened by

f1 (t) f2 (t)
f1 ( )f2 (t )d
(3.18)
Properties of Convolution Integral

Commutative property:
f1 (t) f2 (t) = f2 (t) f1 (t)
(3.19)
Distributive property:
f1 (t) [f2 (t) + f3 (t)] = f1 (t) f2 (t) + f1 (t) f3 (t)
(3.20)
Associate property:
f1 (t) [f2 (t) f2 (t)] = [f1 (t) f2 (t)] f3 (t)
(3.21)
Shift property:
If
f1 (t) f2 (t)
= c(t)
then
f1 (t) f2 (t T )
f1 (t T ) f2 (t)
= c(t T )
= c(t T )
and
f1 (t T1 ) f2 (t T2 ) = c(t T1 T2 )
Convolution with an impulse:

f (t) (t) =
f ( )(t )d = f (t)
(3.22)
(3.23)
Width property: If the durations (widths) of f1 (t) and f2 (t) are T1 and
T2 respectively, the the duration of f1 (t) f2 (t) is T1 + T2 .
92
Zero-State Response and Causality

The zero-state response y(t) of an LTIC system is

f ( )h(t )d
y(t) = f (t) h(t) =
(3.24)
The system is assumed to be linear and time-invariant, there is no other

restrictions on the system and the input signal f (t).
However, in practice, most systems are causal so that their response can
not begin before the input starts. Moreover, most inputs are also causal,
which means they start at t = 0. The causality restrictions on both signals
and systems simplify the limits of integration in equation (3.24). Therefore,
it reduces to
t
f ( )h(t )d
t0
0
y(t) = f (t) h(t) =
(3.25)
0
t<0
Because of commutative property, we can also express (3.25) as
t
h( )f (t )d
t0
0
y(t) =
0
t<0
(3.26)
The results show that if f (t) and h(t) are both causal, the response y(t)
is also causal.
Convolutional Table
The below table lists several pairs of signals and their resulting convolution.
It helps to determine conveniently a system response y(t) to an input f (t),
without performing the integral.
Multiple Inputs
Multiple inputs to LTI systems can be treated by applying superposition
principle. Each input is considered separately, with all other inputs assumed to be zero. When all the inputs are applied simultaneously, the
system output will be the sum of all these individual system response.
Figure 3.2: Convolutional table.
93
94
3.2.2
Graphical Understanding of Convolution
Understanding graphical interpretation of convolution will help:

Evaluating the convolution integral of complicated signals.
To grasp visually/mentally the convolutions result in ltering, sampling and other problems.
When signals have no exact mathematical description, they can be
described only graphically. If the signals are convolved, we should
perform their convolution graphically.
Considering the convolution operation between the signals f (t) and g(t).
If c(t) is the convolution of f (t) with g(t), then

f ( )g(t )d
(3.27)
c(t) = f (t) g(t) =
It is important to remember that this integration is performed with respect

to , t is just a parameter. Therefore, the functions f ( ) and g(t ) should
be sketched as a functions of , not of t.
The graphical explanation of convolution operation is shown in Figure
2.7, Example 2.6, Example 2.7 and Example 2.8 in textbook.
Summary of Graphical Procedure
The procedure for graphical convolution can be summarized as
Keep the function f ( ) xed.
Invert g( ) about the vertical axis ( = 0) to obtain g( ).
Shift the inverted version along the axis by t0 seconds. The shifted
version now represent g(t0 ).
The are under the product of f ( ) and g(t0 ) is c(t0 ), the value of
the convolution at t = t0 .
Repeat this procedure for all value of t to obtain c(t).
Figure 3.3: Graphical explanation of the convolution.
95
96
3.3
Total Response
The total response of a linear system can be expressed as the sum of its
zero-input and zero-state components
n
Total Response =
cj ej t
j=1

zero-input component
f (t) h(t)

zero-state component
(3.28)
Note that for the repeated roots and complex roots the zero-input component should be appropriate modied.
Natural and Forced Responses
The zero-input component is composed of characteristic modes. However,
the zero-state component contains also the characteristic mode terms. This
is generally true for LTI systems. When we lump together all the characteristic mode terms in the total response, it gives us a component known as
the natural response yn (t). The remainder, consisting of noncharacteristic
mode terms, is known as the forced response y (t).
Total Response = Natural Response yn (t) + Natural Response y (t)
(3.29)
3.4
Classical Solution of Dierential Equations
In classical method, to nd the total response of the system, we solve dierential equation to nd the natural component (also known as homogeneous
solution or complementary solution) and forced components (also known
as particular solution), instead of nding the zero-input and the zero-state
components.
3.4.1
Finding Natural Response
The total system response is

y(t) = yn (t) + y (t)
Since y(t) must satisfy the system equation:
Q(D)[yn (t) + y (t) = P (D)f (t)
3.4. CLASSICAL SOLUTION OF DIFFERENTIAL EQUATIONS
97
or
Q(D)yn (t) + Q(D)y (t) = P (D)f (t)
but yn (t) is composed entirely characteristic modes, therefore
Q(D)yn (t) = 0
so that
Q(D)y (t) = P (D)f (t)
(3.30)
The natural response, being a linear combination of the system characteristic modes, has the same form as that of zero-input response; only
its arbitrary constants are dierent. These constants are determined from
auxiliary conditions (systems initial conditions).
3.4.2
Finding Forced Response
It is relatively simple task to nd y (t) of an LTIC system when the input

f (t) has special forms as in Figure 3.4 (also, the Table 2.2, p.141, textbook).
This gure shows a variety of inputs and the form of the forced response
corresponding to each input.
Figure 3.4: Graphical explanation of the convolution.
The coecients in the forced response are determined by substituting

the expression for y (t) in equation (3.30) as
Q(D)y (t) = P (D)f (t)
and then equating coecients of similar terms on both sides of the resulting
expression.
(See example 2.9, pp. 141-142, textbook)
98
Comments on Initial Conditions

In the classical method, the initial conditions are required at t = 0+ . Since
at t = 0 only zero-input component exists. In the classical method, the
zero-input and zero state components can not be separated. Consequently,
the initial conditions must be applied to the total response, which starts at
t = 0+ .
Exponential Input e t
The exponential signal is the most important signal in study LTI system.
The forced response of an exponential input signal has very simple form.
From Figure 3.4 (also, the Table 2.2, p.141, textbook), the forced response
of e t has the form e t. Moreover, it shows that
P ()
(3.31)
= H() =
Q()
Thus, for the input f (t) = e tu(t), the forced response is given by
y (t) = H()e t;
t>0
(3.32)
where H() is given in (3.31).

The total system response y(t) to an exponential input e t is then given
by
y(t) =
n
Kj ej t t + H()e t
(3.33)
j=1
where the arbitrary constants K1 , K2 , , Kn are determined from auxiliary

conditions.
The exponential signal includes a large variety of signals as constant
( = 0), a sinusoid ( = j), an exponentially growing (or decaying)
sinusoid ( = j). We have the forced response of for these cases as
Constant input f (t) = C:
y (t) = CH(0)
(3.34)
Exponential input f (t) = ejt , where = j:

y (t) = H(j)ej
(3.35)
Sinusoidal input f (t) = cos(t + ):

y (t) = |H(j)| cos[t + + H(j)]
(3.36)
3.5. SYSTEM STABILITY
99
Assessment of the Classical Method

The classical method is relative simple than the method of nding zeroinput and zero-state components. However, the classical method has following drawbacks:
The total response obtained by classical method can not be separated
into components arising from internal conditions and the external input. Therefore, we can not express the system response as an explicit
function of input f (t) (this is important in the study of systems).
The classical method is restricted to a certain class of input, it can not
be applied to any input.
In practice, we are mostly to know the initial conditions at t = 0 .
Therefore, we need to derive a new set of auxiliary conditions at t = 0+
from known conditions at t = 0 .
3.5
System Stability
Consider the denition of stability for causal, linear, time-invariant (LTI)

systems. If in the absence of an external input, a system remains in a
particular state indenitely, then that state said to be an equilibrium state
of the system. For an LTI system this equilibrium state is that all initial
conditions are zero. Now if an LTI is in equilibrium, we change this state by
producing nonzero initial conditions. If the system is stable it should return
to equilibrium state or the systems output due to nonzero initial conditions
should approach to zero when t . But the systems output generated
by initial condition (zero-input response) is made up of its characteristic
modes. Therefore, we dene stability as:
A system is (asymptotically) stable if, and only if all its characteristic
modes 0 as t . If any of modes grow without bound as t , the
system is unstable. A system is called marginally stable when the zero-input
response remains bounded (approaches neither zero nor innity), approaching constant or oscillating with a constant amplitude as t .
If an LTI system has n distinct characteristic roots 1 , 1 n , the
zero-input response is given by
n

cj ej t
(3.37)
y0 (t) =
j=1
100 CHAPTER 3. TIME-DOMAIN ANALYSIS OF CONTINUOUS-TIME SYSTEMS
It shows that

lim et =
Re < 0
Re > 0
(3.38)
Consider the location of the systems characteristic in the complex plane.

Assumed that the system has distinct roots only. If a characteristic root
is located in the left half of the complex plane (LHP), its real part is
negative (Re < 0). Similarly, if a characteristic root is located in the
right half of the complex plane (RHP), its real part is positive (Re > 0).
Along the imaginary axis, Re = 0.
According to (3.39), the characteristic modes corresponding to roots in
LHP vanish when t , while the characteristic modes corresponding
to roots in LHP grow without bound when t . The characteristic
modes corresponding to roots on imaginary axis are the form ejt , they are
bounded (neither zero nor innity) as t .
It follows that a system is asymptotically stable if, and only if, all of its
characteristic roots lie on LHP. If any of roots lie in RHP, the system is
unstable. If none of the roots lie in RHP, but if some (unrepeated) roots
lie on imaginary axis, then the system is marginally stable.
Now we consider the modes corresponding to a root repeated r times
e , tet , t2 et , , tr1 et . If repeated roots lie in LHP, the system is stable, since t , tk et 0 (when Re < 0). But if the repeated
roots are on the imaginary axis ( = j), the system is unstable, since
t , tk et .
t
Summary of System Stability

An LTIC is asymptotically stable if, and only if, all of its characteristic
roots lie on LHP. The roots may be unrepeated or repeated.
An LTIC is unstable if, and only if, either one or both of the following
exists: (1) at least one root is in RHP, (2) there are repeated roots on
imaginary axis.
An LTIC is marginally stable if, and only if, there are no roots in RHP,
and there are some unrepeated roots on the imaginary axis.
3.6. THE UNIT IMPULSE RESPONSE H(T )
3.6
101
The Unit Impulse Response h(t)
For a LTIC system specied by the n-order dierential equation

Q(D)y(t) = P (D)f (t)
(3.39)
In practical system, the order n of the polynomial Q(D) and the order m
of the polynomial P (D) are restricted by n m. In most general case,
m = n. Therefore, (3.39) can be expressed as
(Dn +an1 Dn1 + +a1 D+a0 )y(t) = (bn Dn +bn1 Dn1 + +b1 D+b0 )f (t)
(3.40)
In the textbook (pp. 261-263), it shows that the unit impulse response h(t)
is given by
(3.41)
h(t) = bn (t) + [P (D)yn (t)]u(t)
where bn is the coecient of the nth-order term in P (D) in equation (3.40),
yn (t) is a linear combination of the characteristic modes of the system
subject to the following initial conditions:
yn(n1) (0) = 1 and yn (0) = yn (0) = yn (0) = yn(n2) (0) = 0
(3.42)
(k)
where yn (0) is the value of the kth derivative of yn (t) at t = 0. For

example:
n=1
n=2
n=3
n=4
:
:
:
:
yn (0) = 1
yn (0) = 0 and yn (0) = 1
yn (0) = yn (0) = 0 and yn (0) = 1
...
yn (0) = yn (0) = yn (0) = 0 and y n (0) = 1
(3.43)
If m < n, then bn = 0.
Comment
It shows that, the expression of h(t) for all possible values of m and n
is given by
h(t) = P (D)[yn (t)u(t)]
(3.44)
where yn (t) is a linear combination of the characteristic modes of the
system subject to the condition in (3.42).
Other simple method for determination of h(t) is to use the Laplace
transform.
Chapter 4
Signal Representation by Fourier
Series
In this chapter, we consider the issue of representing a signal as a sum of
its components by using Fourier series.
4.1
Signals and Vectors
There is a perfect analogy between signals and vectors. A vector can be

represented as a sum of its components in variety of ways depending on the
choice of coordinate system. A signal can also be represented as a sum of
its components in variety of ways.
4.1.1
Component of a Vector
Let the component of f along x be cx as depicted in Figure 4.1 (that is also

the projection of f on x).
q
cx
Figure 4.1: Component (projection) of a vector along another vector.
102
4.1. SIGNALS AND VECTORS
103
A vector f can be expressed in terms of a certain vector x as

f = cx + e
(4.1)
where e called the error vector. If we approximate f by cx

f cx
(4.2)
then the error in this approximation is the vector e = f cx. We now nd

the unique solution of c about the approximation so that the error vector
is smallest.
Now, the length of the component of f along x is |f | cos , but it is also
c|x| according to Figure 4.1. Therefore,
c|x| = |f | cos
(4.3)
Multiplying both sides by |x| yields

c|x|2 = |f ||x| cos = f .x
where f .x is inner (or dot or scalar) product. Therefore
1
f .x
c=
|x|2
(4.4)
(4.5)
From Figure 4.1, when f and x are perpendicular, or orthogonal, then f

has a zero component along x, consequently, c = 0. From (4.5), we dene
f and x are orthogonal if the inner product of two vectors is zero
f .x = 0
4.1.2
(4.6)
Component of a Signal
Consider the issue of approximating a real signal f (t) in terms of another

real signal x(t) over an interval [t1 , t2 ]:
f (t) x(t) t1 t t2
The error in this approximation is

f (t) cx(t)
e(t) = x(t) =
0
t1 t t2
otherwise
(4.7)
(4.8)
We now select some criterion for the best approximation. We know that
the signal energy is one of possible measure of signal size, which is dened
over the interval [t1 , t2 ] as
t2
t2
Ee =
e2 (t)dt =
[f (t) cx(t)]2 dt
(4.9)
t1
t1
104
CHAPTER 4. SIGNAL REPRESENTATION BY FOURIER SERIES
To minimize, a necessary condition is

dEe
=0
dc
Finally, we obtain
1
c=
=
Ex
(4.10)
t2
f (t)x(t)dt
(4.11)
t1
t
where Ex = t12 x2 (t)dt is energy of x(t). Therefore, the optimum value of
c that minimizes the energy of error signal in the approximation
f (t) cx(t)
is given by equation (4.11).
We dene the real signals f (t) and x(t) are orthogonal if c = 0, or
t2
c=
f (t)x(t)dt = 0
(4.12)
t1
4.1.3
Orthogonality in Complex Signal
Consider the issue of approximating a complex signal f (t) in terms of another complex signal x(t) over an interval [t1 , t2 ]:
f (t) cx(t)
where the coecient c is complex (in general). Similar to the case of real
signals, it shows that (see pp.175-176, textbook)
t2
1
c=
f (t)x (t)dt
(4.13)
Ex t1
t
where Ex = t12 |x(t)|2 dt is energy of complex signal x(t). Therefore, if two
complex signals f (t) and x(t) are orthogonal, then
t2
t2
f (t)x (t)dt = 0 or
f (t)x(t)dt = 0
(4.14)
t1
4.2
4.2.1
t1
Signal Comparison: Correlation

Correlation Coecient
In this section we nd the similarity measure or the correlation coecient

cn of two signals.
4.2. SIGNAL COMPARISON: CORRELATION
105
We consider the signals over time interval from to . We dene

the correlation coecient cn of two real signals f (t) and x(t) as

1
c=
f (t)x(t)dt
(4.15)
Ef Ex
where Ef and Ex are energies of f (t) and x(t), respectively. Note that
cn is independent of the energies of f (t) and x(t). From (4.15), using the
Schwarz inequality, we obtain
1 cn 1
(4.16)
When cn = 0, the similarity is zero, f (t) and x(t) are orthogonal.

Therefore, orthogonal signals are unrelated signals.
When cn = 1, the similarity is maximum.
When cn = 1, the dissimilarity is maximum.
We can also dene the correlation coecient cn for two complex signals
f (t) and x(t) as

1
f (t)x (t)dt
(4.17)
c=
Ef Ex
(See application of correlation coecient in signal detection on pp. 179181, textbook)
4.2.2
Correlation Functions
When f (t) and g(t) are similar but they are time-shifted so that nonoverlapping in time. From (4.15), we obtain

1
f (t)x(t)dt = 0
(4.18)
cn =
Ef Ex
Thus, the correlation is zero because two similar signals are disjoint. For
this reason, instead of using the integral on the right-hand of (4.15), we
dene the crosscorrelation function of two real signals f (t) and g(t) as

f g (t) =
f ( )g( t)d
(4.19)
where is a dummy variable, and the signal g( t) is the signal g( )

delayed by t seconds with respect to the signal f ( ). Therefore, f g (t)
106
is an indication of similarity (correlation) of the signal f with signal g

delayed by t seconds. Thus, the crosscorrelation function f g (t) measures
the similarity of signals even if they are disjoint.
For the complex signals, we dene

f ( )g( t)d
(4.20)
f g (t) =
Correlation Function and Convolution

In correlation in (4.19), g( t) is the g( ) time-shifted by t (without timeinversion). In convolution, we follow the same procedure, except that g is
time-inverted before it is shifted by t. Therefore,
f g (t) = f (t) g(t)
(4.21)
Autocorrelation Function
Correlation of a signal with itself is called autocorrelation. The autocorrelation function f (t) of a signal f (t) is dened as

f (t) =
f ( )f ( t)d
(4.22)
4.3
Signal Representation by Orthogonal Signal Set
Real Signals
We dene orthogonality of real signal set x1 (t), x2 (t), , xN (t) over interval
[t1 , t2 ] as

0
m = n
xm (t)xn (t)dt =
(4.23)
En m = n
If En = 1 for all n, then the set is normalized and is called an orthonormal

set.
Now we consider approximation a signal f (t) over interval [t1 , t2 ] by a
set of N real, mutually orthogonal signal x1 (t), x2 (t), , xN (t) as
f (t) c1 x1 (t) + c2 x2 (t) + + cN xN (t) =
N

n=1
cn xn (t)
(4.24)
4.3. SIGNAL REPRESENTATION BY ORTHOGONAL SIGNAL SET
107
The error e(t) in the approximation (4.24) is

e(t) = f (t)
N
cn xn (t)
(4.25)
n=1
It is shown in the Appendix 3A (textbook) that Ee , the energy of the

error signal e(t), is minimized if we choose
t2
t2
1
t1 f (t)xn (t)dt
c n = t2
f (t)xn (t)dt
n = 1, 2, , N (4.26)
=
2
E
n
t1
t1 xn (t)dt
For this choice of cn , the error signal energy Ee is given by
t2
N

2
Ee =
f (t)dt
c2n En
t1
(4.27)
n=1
If N , then the error energy 0. In this case the orthogonal signal

set is said complete, and (4.24) is no more approximation but an equality
on t1 t t2
f (t) = c1 x1 (t) + c2 x2 (t) + + cn xn (t) + =
cn xn (t)
(4.28)
n=1
where the coecients cn are given by (4.26). The series on the right-hand
side of (4.28) is called generalized Fourier series of f (t) with respect to
the set {xn (t)}. The set {xn (t)} is called a set of basis functions or basis
signals.
It follow also that the energy of f (t) is now equal to the sum of the
energies of its orthogonal components c1 x1 (t), c2 x2 (t), . Thus
t2

2
2
2
f (t)dt = c1 E1 + c2 E2 + =
c2n En
(4.29)
t1
n=1
It is Parsevals theorem.
Complex Signals
The above results can be generalized to complex signal as follows: A set
of functions x1 (t), x2 (t), , xN (t) is mutually orthogonal over the interval
[t1 , t2 ] if

t2
0
m = n
xm (t)xn (t)dt =
(4.30)
En m = n
t1
108
If this set is complete for a certain class of functions, then a function

f (t) in this class can be expressed as
f (t) = c1 x1 (t) + c2 x2 (t) + + cn xn (t) +
(4.31)
t2
1
cn =
f (t)xn (t)dt
(4.32)
En t1
An interesting property of the coecients c1 , . . . , cn , . . . cN is the optimum value of any coecient is independent of the number of terms used in
the approximation of a signal f (t). It means we can continue to add terms
to the approximation without disturbing the previous terms.
where
4.4
Trigonometric Fourier Series
Consider a signal set

{1, cos 0 t, cos 20 t, . . . , cos n0 t, . . . ; sin 0 t, sin 20 t, . . . , sin n0 t, . . .}
(4.33)
A sinusoid of frequency n0 t is called the nth harmonic of the sinusoid
of frequency 0 t when n is an integer. In this set sinusoid of frequency 0 t
is called the fundamental frequency. In Appendix 3B, it shows that this
signal set is orthogonal over any interval T0 = 2/0 , where T0 is the period
of the fundamental frequency. It has shown that

t1 +T0
0
m = n
(4.34)
(cos n0 t)(cos m0 t)dt = T0
m = n = 0
t1
2

t1 +T0
0
m = n
(sin n0 t)(sin m0 t)dt = T0
(4.35)
m
=
n
=

0
t1
2
and
t1 +T0
(sin n0 t)(sin m0 t)dt = 0
for all n and m
(4.36)
t1
where t1 is arbitrary value.

These equation show that the signal set (4.33) is orthogonal over any
interval of duration T0 . This is the trigonometric set and shown to be a
complete set. Therefore, we can express a signal f (t) by a trigonometric
Fourier series over any duration T0 as
f (t) = a0 + a1 cos 0 t(t) + a2 cos 20 t(t) +
t 1 t t 1 + T0
+b1 sin 0 t(t) + b2 sin 20 t(t) +
(4.37)
4.4. TRIGONOMETRIC FOURIER SERIES
109
or
f (t) = a0 +
(an cos n0 t + bn sin n0 t)
t1 t t1 + T0
(4.38)
n=1
where 0 = 2
T0 .
Using (4.26) and (4.34)-(4.36), we can determine the Fourier coecients
a0 , an and bn as

1 t1 +T0
a0 =
f (t)dt
(4.39)
T0 t 1

2 t1 +T0
f (t) cos n0 tdt
n = 1, 2, 3, . . .
(4.40)
an =
T0 t 1

2 t1 +T0
bn =
f (t) sin n0 tdt
n = 1, 2, 3, . . .
(4.41)
T0 t 1
Compact Trigonometric Fourier Series
The trigonometric Fourier series in (4.38) contains sine and cosine terms of
the same frequency, then we can express in the compact form the trigonometric Fourier series as

f (t) = C0 +
Cn cos(n0 t + n )
t1 t t1 + T0
(4.42)
n=1
where the coecient Cn and n are computed from an and bn as

C0 = a
0
Cn =
a2n +b2n
bn
n = arctan
an
(4.43)

Periodicity of the Trigonometric Fourier Series
It is shown in the textbook that an arbitrary signal f (t) may be expressed
as a trigonometric Fourier series over any interval of T0 seconds. Outside
this interval, what happens to the Fourier series.
It is shown that the trigonometric Fourier series is a periodic function
of period T0 (the period of the fundamental frequency). Let denote the
trigonometric Fourier series on the right-hand side of (4.42) by (t), then
(t + T0 ) = (t)
for all t
(4.44)
110
Thus, when we present a signal f (t) by trigonometric Fourier series over a

certain duration T0 , the function f (t) and its Fourier series (t) need be
equal only over that interval of T0 . Outside this interval, the Fourier series
(t) repeat periodically with period T0 .
Now, if the function f (t) were itself to be periodic with period T0 , then
a Fourier series representing f (t) over an interval of T0 will also represent
f (t) for all t (not just over the interval T0 ).
Another interesting fact is that a periodic signal f (t) can be generated
by a periodic repetition of any of its segment of duration T0 (see Figure 1.7,
p. 60, textbook). Therefore, the trigonometric Fourier series representing
a segment of f (t) of duration T0 starting at any instant represents f (t) for
all t. It follows that in computing the coecients a0 , an and bn , we may use
any value for t1 in (4.39)-(4.41).
The Fourier Spectrum

The compact trigonometric Fourier series in (4.42) indicates that a periodic
signal f (t) can be expressed as a sum of sinusoids frequency:
0, 0 , 20 , . . . , n0 , . . . ,
whose amplitudes are
C0 , C1 , . . . , Cn , . . . ,
and whose phase are
0, 1 , 2 , . . . , n , . . . ,
respectively.
We now can plot amplitude Cn vs. (amplitude spectrum) and n vs.
(phase spectrum). These plots together are the frequency spectra of f (t).
The spectra show amplitudes and phases of various sinusoidal components of (t). Knowing frequency spectra, we can reconstruct or synthesize
(t). Therefore, the frequency spectra provide an alternative description
of signal: the frequency-domain description. A signal, therefore, has a dual
identity: the time-domain identity (t) and the frequency-domain identity
(Fourier spectra). The two identities complement each other providing a
better understanding of a signal.
4.4. TRIGONOMETRIC FOURIER SERIES
4.4.1
111
The Eect of Symmetry
If f (t) is an even periodic function of t, the Fourier coecients can be

computed as

2 T0 /2
a0 =
f (t)dt
(4.45)
T0 0

4 T0 /2
f (t) cos n0 tdt
(4.46)
an =
T0 0
bn = 0
(4.47)
If f (t) is an odd periodic function of t, the Fourier coecients can be
computed as
a0 = an = 0
(4.48)
T0 /2
4
bn =
f (t) sin n0 tdt
(4.49)
T0 0
4.4.2
Determining the Fundamental Frequency and Period
Consider a sum of sinusoids of any frequencies. We need to check whether

it represents a periodic signal. If so, how to determine the period?
Recall that every frequency in a periodic signal is an integral multiple of
the fundamental frequency 0 . Therefore, the ratio of any two frequencies
is a rational number m/n, where m and n are integers. In this case, two
frequency are said to be harmonically related.
The largest positive number of which all the frequencies are integral
multiples is the fundamental frequency.
For example, consider the following functions
1
2
7
f1 (t) = 2 + 7 cos( t + 1 ) + 7 cos( t + 2 ) + 5 cos( t + 3 )
2
3
6
and
f2 (t) = 2 cos(2t + 1 ) + 5 sin(t + 2 )
The frequencies in the spectrum of f1 (t) are 12 , 23 and 76 . The ratios of the
successive frequencies are 34 and 47 , respectively. Because these numbers are
rational, all three frequencies are harmonically related and therefore signal
f1 (t) is periodic. The largest number of which 12 , 23 and 76 are integral multiples is 16 . Therefore, the fundamental frequency is 16 . The three frequencies
in the spectrum are the third, fourth, and seventh harmonics.
The signal f2 (t) is not periodic because the ratio of two frequencies in
the spectrum is 2 , which is not a rational number.
112
4.4.3
Role of Amplitude and Phase Spectra in Wave Shaping
Student is required to read the Section 3.4-3, pp. 201-205, textbook, to

understand:
How to synthesize any signal f (t) by using proper combination of
amplitudes and phases of various sinusoids.
The Gibbs phenomenon in Fourier synthesis of discontinuous functions
(e.g. square function).
Using Matlab for demonstration is highly recommended.
4.5
Exponential Fourier Series
It is shown in Appendix 3C in the textbook that the set of exponentials:

ejn0 t ,
n = 0, 1, 2, . . .
is orthogonal over any interval of duration T0 = 2

0 and is a complete set.
Therefore, from (4.31) and (4.32), a signal f (t) can be expressed over an
interval T0 as an exponential Fourier series:
f (t) =
Dn ejn0 t
(4.50)
f (t)ejn0 t dt
(4.51)
n=
where
1
Dn =
T0
t1 +T0
t1
for any value of t1 . The exponential Fourier series in (4.50) is periodic

with period T0 . This is basically another form of the trigonometric Fourier
series.
The connection between the coecients in the trigonometric Fourier
series (4.42) and the exponential Fourier series (4.50) is
Dn =
Dn =
or
1
jn
2 Cn e
1
jn
2 Cn e
1
Dn = (an jbn )
2
(See example 3.6, pp. 208, textbook)
(4.52)
(4.53)
4.5. EXPONENTIAL FOURIER SERIES
4.5.1
113
Exponential Fourier Spectra
For the exponential spectra, we plot the magnitude |Dn | and the angle Dn
of Dn as functions of .
Compare (4.39) and (4.51) for n = 0 shows that
D0 = a0 = C0
(4.54)
Equation (4.52) shows that, for real f (t)

1
|Dn | = |Dn | = Cn
2

Dn = n
and
Thus,
Dn = |Dn |ejn
n = 0
Dn = n
and Dn = |Dn |ejn
(4.55)
(4.56)
(4.57)
It follows that the magnitude spectrum (|Dn | vs. ) is an even function

of and the angle spectrum ( Dn vs. ) is an odd function of , when
f (t) is a real signal.
Negative Frequency in Exponential Fourier Spectra?
Student is required to explain what is a negative frequency in the exponential Fourier spectra. Read pp. 209-211 in the textbook.
Bandwidth of a Signal
The dierence between the highest and the lowest frequencies of the (trigonometric) spectra components of a signal is the bandwidth of the signal.
4.5.2
Parsevals Theorem
The power of periodic signal f (t) in terms of its trigonometric Fourier series
is given by
1 2
2
C
(4.58)
Pf = C0 +
2 n=1 n
It is also shown that the power of periodic signal f (t) in terms of its exponential Fourier series is given by
Pf =

n=
|Dn |2
(4.59)
114
For real f (t), |Dn | = |Dn |. Therefore

Pf =
D02
+2
|Dn |2
(4.60)
n=1
(See example 3.9, pp. 215, textbook. Study how to compute harmonic
distortion).
Limitations of the Fourier Series Method of Analysis
We have studied in this Chapter a method of representing a period signal
as a weighted sum of everlasting exponentials whose frequency lie along
the j-axis in the s-plane. This Fourier series representation is valuable in
many applications. However, as a tool for analyzing linear systems, it has
some limitations as
The Fourier series can be used only for periodic signal. However, practical signals are aperiodic. This limitation can be overcome by representing aperiodic signals in terms of everlasting exponentials through
the Fourier integral (which may be considered as an extension of the
Fourier series). We will consider the Fourier integral in Chapter 5.
The Fourier series technique can be applied for stable systems. It
can not handle easily unstable or marginally stable systems. This
limitation can be overcome by using Laplace integral, which will be
discussed in Chapter 6.
Chapter 5
Continuous-Time Signal Analysis:
The Fourier Transform
In previous chapter, we considered the issue of representing periodic signals as a sum of sinusoids or exponentials. In this chapter we extend this
representation to aperiodic signals.
5.1
Aperiodic Signal Representation by Fourier Integral
We now study how an aperiodic signal can be expressed as a continuous

sum (integral) of exponentials.
To represent an aperiodic signal f(t) such as the one depicted in Figure 5.1(a) by exponential signals, we construct a new periodic signal fT0 (t)
formed by repeating the signal f (t) at intervals of T0 seconds, as in Figure 5.1(b). The period T0 is made long enough to avoid overlap between
the repeating pulses.
This periodic signal fT0 (t) can be represented by an exponential Fourier
series. If we let T0 , the pulses in the periodic signal repeat after an
innite interval and therefore
lim fT0 (t) = f (t)
T0
(5.1)
Thus, the Fourier series representing fT0 (t) will also represent f (t) in the
limit T0 . The exponential Fourier series for fT0 (t) is given by
fT0 (t) =

n=
115
Dn ejn0 t
(5.2)
116CHAPTER 5. CONTINUOUS-TIME SIGNAL ANALYSIS: THE FOURIER TRANSFORM
f(t)
0
(a)
fT0(t)
-T 0
T0
(b)
Figure 5.1: Construction of periodic signal fT0 (t) by periodic extension of aperiodic signal f (t).
5.1. APERIODIC SIGNAL REPRESENTATION BY FOURIER INTEGRAL
where
1
Dn =
T0
T0 /2
T0 /2
fT0 (t)ejn0 t dt
117
(5.3)
and
2
(5.4)
T0
Observe that integrating fT0 (t) over (T0 /2, T0 /2) is the same as integrating f (t) over (, ). Therefore, (5.3) can be expressed as

1
f (t)ejn0 t dt
(5.5)
Dn =
T0
0 =
We dene F (), a continuous function of , as

F () =
f (t)ejt dt
(5.6)
From (5.5) and (5.6), it shows that

Dn =
1
F (n0 )
T0
(5.7)
It means that the Fourier coecients Dn are (1/T0 ) times the samples of
F () uniformly spaced at interval of 0 . Therefore, (1/T0 )F () is the
envelope for the coecients Dn .
Substitution of (5.7) in (5.2) yields

F (n0 ) jn0 t
e
fT0 (t) =
T
0
n=
(5.8)
As T0 , 0 0 (innitesimal). Hence, we replace 0 by a notation

, then (5.4) becomes
2
(5.9)
=
T0
and (5.8) becomes

F (n) (jn)t
e
(5.10)
fT0 (t) =
2
n=
This equation shows that fT0 (t) can be expressed as a sum of exponentials
of frequencies 0, , 2, . . .. In the limit as T0 , 0 and
fT0 (t) f (t). Therefore
1
f (t) = lim fT0 (t) = lim
F (n)e(jn)t
T0
0 2
n=
(5.11)
The sum on the right-hand side of (5.11) can be seen as the area under
the function F ()e(j)t . Therefore

1
F ()ejt d
(5.12)
f (t) =
2
The integral on the right-hand side is called the Fourier integral. From
(5.12), it means that an aperiodic signal f (t) is represented by a Fourier
integral (rather then Fourier series). This integral is basically a Fourier
series (in the limit) with fundamental frequency 0. The amount of
the exponential ejnt is F (n)/2. Thus, the function F () given
by (5.6) acts as a spectral function.
We call F () the direct Fourier transform of f (t) and f (t) the inverse
Fourier transform of F (). We call f (t) and F () are a Fourier transform
pair. Symbolically, this statement is expressed as
F () = F[f (t)] and f (t) = F 1 [F ()]
or
f (t) F ()

where
F () =
f (t)ejt dt
(5.13)
and
1
f (t) =
2
F ()ejt d
(5.14)
We can plot F () as a function of . Since F () is complex, we have

both amplitude and phase spectra
F () = |F ()|ej
F ()
(5.15)
where |F ()| is the amplitude and F () is the phase of F (). From (5.13),
we have

f (t)ejt dt
(5.16)
F () =
From (5.13) and (5.16), it follows that if f (t) is a real function of t, then
F () and F () are complex conjugates. Therefore
|F ()| = |F ()|
F () = F ()
(5.17)
5.1. APERIODIC SIGNAL REPRESENTATION BY FOURIER INTEGRAL
119
Thus, for real f (t), the amplitude spectrum |F ()| is an even function, and
the phase spectrum is an odd function of . The transform F () is the
frequency-domain specication of f (t).
Note that Fourier spectrum of a periodic signal has nite amplitude
and exists at discrete frequency (fundamental frequency and its multiples),
while the spectrum of an aperiodic signal continuous.
LTIC System Response Using the Fourier Transform

If h(t) is the systems unit impulse response, then the system response y(t)
to an input exponential est is given by

st
s(t )
st
y(t) = h(t) e =
h( )e
d = e
h( )es d
(5.18)
The integral in the right-hand side is a function of s. We dene the

system transfer function H(s) as

H(s) =
h(t)est dt
(5.19)
Thus,
y(t) = H(s)est
Setting s = j in this equation yields
H(j) = H() =
h(t)ejt dt
(5.20)
(5.21)
The right-hand side is the Fourier transform of h(t).

We represented a signal f (t) as a sum of exponentials, thus we can
nd a system response to f (t) as a sum of the systems responses to the
exponential components of f (t).
Consider a asymptotically stable LTIC system with transfer function
H(s). The response of this system to exponential ejt is H()ejt .
It shows in the textbook that

1
1
jt
F ()H()e d =
Y ()ejt d
(5.22)
y(t) =
2
2
where Y (), the Fourier transform of y(t), is given by
Y () = F ()H()
(5.23)
Thus, transmission of a signal through a linear system is seen as transmission of various sinusoidal components of the signal through the system.
5.2
Transforms of Some Useful Functions

f (t)
F ()
eat u(t)
1
,
a + j
a>0
eat u(t)
1
,
a j
a>0
ea|t|
2a
,
a2 + 2
a>0
teat u(t)
1
,
(a + j)2
tn eat u(t)
n!
,
(a + j)n+1
(t)
2()
ej0 t
2( 0 )
cos 0 t
[( 0 ) + ( + 0 )]
a>0
a>0
10 sin 0 t
j[( + 0 ) ( 0 )]
11 u(t)
() +
12 sgnt
2
j
13 cos 0 tu(t)
j
[( 0 ) + ( + 0 )] + 2
2
0 2
14 sin 0 tu(t)
0
[( 0 ) ( + 0 )] + 2
2j
0 2
15 eat sin 0 tu(t)
0
,
(a + j)2 + 02
1
j
a>0
5.2. TRANSFORMS OF SOME USEFUL FUNCTIONS
16 eat cos 0 tu(t)

t
17 rect
sinc
21
Wt
2
/2 2

2

2
sinc
2
4
(t nT )
n=
22 et
a>0

rect
2W
W
sinc(W t)
18

t
19
W
20
sinc2
2
a + j
,
(a + j)2 + 02
121

2W
0
( n0 ),
n=
2 2
2e /2
0 =
0
T
5.3
Some Properties of Fourier Transform
Operation
f (t)
F ()
Addition
f1 (t) + f2 (t)
F1 () + F2 ()
Scalar multiplication
kf (t)
kF ()
Symmetry
F (t)
2F ()
Scaling (a is real)
f (at)
1
F
|a|
a
Time shift
f (t t0 )
F ()ejt0
Frequency shift (0 real)
f (t)ejt0
F ( 0 )
Time convolution
f1 (t) f2 (t)
F1 ()F2 ()
Frequency convolution
f1 (t)f2 (t)
1
F1 () F2 ()
2
Time dierentiation
dn f
dtn

(j)n F ()
f (x)dx
Time integration
5.4
F ()
+ F (0)()
j
Signal Transmission through LTIC Systems
If f (t) and y(t) are the input and output of an LTIC system with transfer
function H(), then
Y () = H()F ()
(5.24)
This result applies only to asymptotic (and marginally) stable systems.

Moreover, f (t) has to be Fourier transformable. (See example 4.15, p. 267,
textbook).
5.4. SIGNAL TRANSMISSION THROUGH LTIC SYSTEMS
123
Signal Distortion during Transmission

From (5.24), we obtain
|Y ()| = |H()F ()|
Y () = F () + H()
(5.25)
We can see that during transmission, the input amplitude spectrum

|F ()| is changed to |F ()||H()|. Similarly, the input phase spectrum
F () is changed to F () H(). Clearly, |H()| is the amplitude response and H() is the phase response of the system. The H() is called
the frequency response of the system.
In several applications, such as signal amplication or message signal
transmission over a communication channel, it is required that the output
signal is a replica of the input signal. In such cases, we need to minimize the
distortion caused by the amplier or the communication channel. Therefore, it is practical interest to determine the characteristics of the system
that allows a signal to pass without distortion (distortionless transmission).
In distortionless transmission, the input f (t) and output y(t) satisfy the
condition
y(t) = kf (t td )
(5.26)
where td is the delay of the output with respect to the input, k is a constant.
The Fourier transform of this equation is
Y () = kF ()ejtd
(5.27)
Y () = F ()H()
(5.28)
H() = kejtd
(5.29)
But
Therefore
This is the required transfer function of the system for distortionless transmission. It follows that
|H()| = k
(5.30)
H() = td
This result shows that for distortionless transmission, the amplitude response |H()| must be a constant, and the phase response H() must be
a linear function of .
5.5
Ideal and Practical Filters
Ideal lters allow distortionless transmission of a certain band of frequencies

and suppress all the remaining frequencies. For example, the ideal low pass
lter allows all components below = W rad/s to pass without distortion
and suppressed all components above = W . If the ideal low pass lter
has a linear phase of td , which results in a time delay of td seconds for
all its input components of frequencies below W rad/s. Therefore, if the
input signal is f (t) bandlimited to W rad/s, the output y(t) is f (t) delayed
by td as
y(t) = f (t td )
The ideal low pass lter in this case has the amplitude response

|H()| = rect
2W
and the phase response

H() = ejtd
The frequency response of the ideal low pass lter is then

ejtd
H() = rect
2W
The unit impulse response h(t) of this lter is obtained as

W
1
jtd
e
sinc[W (t td )]
rect
=
h(t) = F
2W
(5.31)
(5.32)
The response h(t) begins even before input applied at t = 0, then this ideal
low pass lter is noncausal and therefore physically unrealizable.
In the frequency domain, this condition equivalent to the well-known
Paley-Wiener criterion, which states that the necessary and sucient condition for the amplitude response |H()| to be realizable is

| ln |H()||
d <
(5.33)
1 + 2
If |H()| = 0 over any nite band, | ln |H()| = over that band, and the
condition (5.33) is violated. If, however, H() = 0 at a single frequency
(or a set of discrete frequencies), the condition in (5.33) may still nite.
According to this criterion, ideal lters (low pass, high pass, band pass
lters) are unrealizable.
5.6. DATA TRUNCATION: WINDOW FUNCTIONS
125
The unit impulse response h(t) in (5.32) is not realizable. One approach
is to cut o the tail of h(t) for t < 0. The resulting causal impulse response
h(t)
is
= h(t)u(t)
h(t)
The truncation operation, however, creates some unsuspected problems as
discussed in next section.
In practice, we can realize a variety of lter characteristics to approach
ideal characteristics. Practical (realizable) lter characteristics are gradual,
without jump discontinuities in amplitude response. We will study them
in Chapter 7.
5.6
Data Truncation: Window Functions
We often need to truncate data in diverse situation from numerical computation to lter design. Data truncation can occur in both time and
frequency domain.
Window Functions
Truncation operation may be regarded as multiplication a signal of a large
width by a window function of a smaller nite width.
Consider a signal f (t) and a window function w(t). If f (t) F ()
and w(t) W (), and if the windowed function fw (t) Fw (), then
fw (t) = f (t)w(t)
and
1
F () W ()
2
Truncation operation using window function causes two eects:
Fw () =
1. The spectral spread or wider mainlobe width of the windowed signal

Fw (). This can be explained by using the width property of convolution. We know that the signal bandwidth is inversely proportional to
the signal width. Therefore, to reduce the spectral spread, we need to
increase the window width.
2. In fact, the spectrum of W () is really not strictly bandlimited. This
causes the spectrum of Fw () is also not strictly bandlimited, even
though the F () may be strictly bandlimited. This eect is called
leakage or sidelobe. We can achieve a given leakage behavior by selecting a suitably smooth window (e.g. Hanning or Hamming windows).
For a given window width, the remedies for the two eects are incompatible. For example, among all the windows of a given width, the rectangular
window has the smallest spectral spread, but has high level sidelobes. which
decay slowly. A smooth window of the same width has smaller and faster
decaying sidelobes, but it has wider mainlobe. However, we can remedy
both the side eects of truncation by selecting a suitably smooth window
of sucient width.
There are several well-known window functions, such as Bartlett (triangular), Hanning, Hamming, Blackman and Kaiser, which truncate the data
gradually. These window oer dierent tradeos with respect to spectral
spread (A), the peak sidelobe magnitude (C), and the leakage rollo rate
(B) as indicated below.
Window
w(t)
Rectangular

t
rect
T
4
T
13.3
8
T
12
26.5
8
T
18
31.5
8
T
42.7
8
T
18
58.1
11.2
T
59.9
Bartlett
t

2T
Hanning

2t
0.5 1 + cos
T
Hamming
2t
0.54 + 0.46 cos
T

2t
+
0.42 + 0.5 cos

T
4t
+0.08 cos
T
Blackman
Kaiser

I0 1 4(t/T )2
I0 ()
1 10
B
[dB/oct]
C
[dB]
( = 8.168)
5.7. SIGNAL ENERGY
5.7
127
Signal Energy
In the Chapter 2, the signal energy Ef of a signal f (t) was dened as

Ef =
|f (t)|2 dt
(5.34)
The well-known Parsevals theorem for Fourier transform is given by (see

detail in the textbook)

1
2
Ef =
|f (t)| dt =
|F ()|2 d
(5.35)
2
where F () is the Fourier transform of f (t). The |F ()|2 is the energy

spectral density.
Essential Bandwidth of a Signal
Spectral of most of the signals extend to innity. However, because the
energy of any practical signal is nite, the signal spectrum must approach
0 as . Most of the signal energy is contained within a certain band
of B Hz, and the energy contributed by the components beyond B Hz is
negligible. We can therefore suppress the signal spectrum beyond B with
little eect on the signal shape and energy. The bandwidth B is called the
essential bandwidth of the signal. The criterion for selecting B depends on
the error tolerance in a particular application.
Energy Spectral Density and Autocorrelation Function
The autocorrelation function f (t) for a real f (t) with itself is given by

f (t) =
f (x)f (x t)dx = f (t) f (t)
(5.36)
It follows,
f (t) = f (t)
Therefore, from (5.36), the Fourier transform of the autocorrelation is its
energy spectral density |F ()|2
f (t) |F ()|2
(5.37)
Chapter 6
Continuous-Time System Analysis:
The Laplace Transform
The Fourier transform allows us to represent a signal f (t) as a continuous
sum of exponentials of the form ejt , whose frequencies are restricted to
the imaginary axis (s = j). This representation is valuable in the analysis
and processing of signals.
In system analysis, however, the Fourier transform has some disadvantages. First, the Fourier transform exists only for a restricted class of
signals, and therefore, can not be used for such inputs (such as growing
exponential signals). Second, the Fourier transform can not be used easily
to analyze unstable or marginally stable systems.
6.1
The Laplace Transform
The basic reason for both diculties is that for some signals. such as
et u(t), ( > 0), the Fourier does not exist because the ordinary sinusoids
or exponentials of the form ejt are incapable of synthesizing exponentially
growing signals.
This problem can be resolved if we can use the basis signal of the form est
(instead of ejt ). This is known as the bilateral Laplace transform, where the
frequency variable s = j is generalized to s = + j. Such generalization
permits us to use exponentially growing sinusoids to synthesize a signal
f (t).
128
6.1. THE LAPLACE TRANSFORM
6.1.1
129
Analytical Development of the Bilateral Laplace Transform
The Fourier transform is dened by

F (j) =
f (t)ejt dt
(6.1)

1
f (t) =
F (j)ejt d
2
Consider now the Fourier transform of f (t)et ( real)

t
t jt
f (t)e e
dt =
f (t)e(+j)t dt
F[f (t)e ] =
and
(6.2)
(6.3)
It follows from (6.1) that the above integral is F ( + j). Thus

F[f (t)et ] = F ( + j)
(6.4)
The inverse Fourier transform of the above equation yields

1
F ( + j)ejt d
(6.5)
f (t)et =
2
then

1
f (t) =
F ( + j)e(+j)t d
(6.6)
2
Now change the variable from to s = + j. The limits of integration
become from (c j) to (c + j), where c > 0 . 0 is a certain minimum
value of for a given f (t) so that the Laplace transform exists. The region
in the complex plane where > 0 is called region of convergence. Thus,
(6.6) becomes

1
F (s)est ds
(6.7)
f (t) =
2j
From (6.3) and (6.4), we obtain

f (t)est dt
(6.8)
F (s) =
The pair of equations (6.7) and (6.8) is known as bilateral Laplace transform
pair or two-sided Laplace transform pair. The bilateral Laplace transform
is denoted symbolically as
F (s) = L[f (t)] and f (t) = L1 [F (s)]
or
f (t) F (s)
130CHAPTER 6. CONTINUOUS-TIME SYSTEM ANALYSIS: THE LAPLACE TRANSFORM
Response of LTIC Systems

Similar to the Fourier transform, it is shown that for an LTIC system
with transfer function H(s), if the input and the output are f (t) and y(t),
respectively, and if
f (t) F (s),
y(t) Y (s)
then
Y (s) = F (s)H(s)
(6.9)
Region of Convergence
The region of convergence for F (s) is the set of values of s for which the
integral in (6.8) converges. See example 6.1, pp. 366, textbook.
Unilateral Laplace Transform
The unilateral Laplace transform is a special case of the bilateral Laplace
transform, where all signals are restricted to being causal. The unilateral
Laplace transform is dened as

F (s)
f (t)est dt
(6.10)
0
We will consider the unilateral Laplace transform and its application to

system analysis.
Existence of the Laplace Transform
The existence of the Laplace transform is guaranteed if

|f (t)et |dt <
(6.11)
For the unilateral Laplace transform, there is a unique inverse transform

of F (s), therefore, there is no need to specify the region of convergence
explicitly.
6.1.2
Finding the Inverse Transform
We can nd the inverse transform from the below table. Most of the transform F (s) of practical interest are rational functions, that is, ratio of polynomials in s. Such functions can be expressed as a sum of simpler functions
of the form listed in the table by using partial fraction expansion.
6.1. THE LAPLACE TRANSFORM
131
f (t)
F ()
(t)
u(t)
1
s
tu(t)
1
s2
tn u(t)
et u(t)
1
s
tet u(t)
1
(s )2
tn et u(t)
n!
(s )n+1
8a
cos btu(t)
s
s2 + b 2
8b
sin btu(t)
b
s2 + b 2
9a
eat cos btu(t)
s+a
(s + a)2 + b2
9b
eat sin btu(t)
b
(s + a)2 + b2
10a
reat cos(bt + )u(t)
(r cos )s + (ar cos br sin )

s2 + 2as + (a2 + b2 )
10b
re
at
cos(bt + )u(t)
n!
sn+1
0.5rej
0.5rej
+
s + a jb s + a + jb
10c
reat cos(bt + )u(t)
As + B
s2 + 2as + c
r=
=
b=
A2 c+B 2 2ABa
,
ca
2

AaB
arctan Aca2 ,
c a2

10d
eat
b=

B Aa
A cos bt +
sin bt u(t)
b
c a2
As + B
s2 + 2as + c
6.2. PROPERTIES OF THE LAPLACE TRANSFORM
6.2
133
Properties of the Laplace Transform
Operation
f (t)
F ()
Addition
f1 (t) + f2 (t)
F1 (s) + F2 (s)
Scalar multiplication
kf (t)
kF (s)
Scaling
f (at), a 0
1 s
F
a
a
Time shift
f (t t0 )u(t t0 )
F (s)est0 , t0 0
Frequency shift
f (t)es0 t
F (s s0 )
Time convolution
f1 (t) f2 (t)
F1 (s)F2 (s)
Frequency convolution
f1 (t)f2 (t)
1
F1 (s) F2 (s)
2
Time dierentiation
df
dt
sF (s) f (0 )
d2 f
dt2
s2 F (s) sf (0 ) f(0 )
d3 f
dt3

s3 F (s) s2 f (0 )
sf(0 ) f(0 )
Time integration
f ( )d
1
F (s)
s
f ( )d
1
1
F (s) +
s
s
Frequency dierentiation
tf (t)
dF (s)
ds
f (t)dt
Frequency integration
f (t)
t
Initial value
f (0+ )
Final value
f ()
F (z)dz
s
lim sF (s) (n > m)
lim sF (s)
(poles of sF (s) in LHP)

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Analog Signal Processing

Dr.-Ing. Tuan Do-Hong

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The quantities v and i are generally function of time. Therefore P is a

1.1. CIRCUIT INTRODUCTORY

A circuit consists of circuit elements attached to each other with ideal

CHAPTER 1. FUNDAMENTALS OF CIRCUIT ANALYSIS

Kirchho Current Law (KCL): KCL follows from our assumption

Figure 1.5: Example for KCL.

Kirchho Voltage Law (KVL): KVL follows from the denition of

1.1. CIRCUIT INTRODUCTORY

Figure 1.6: Example for KVL.

Linear Circuit Elements

CHAPTER 1. FUNDAMENTALS OF CIRCUIT ANALYSIS

P > 0: Resistor always absorbs power.

Independent Voltage Source (IVS)

Figure 1.8: Independent voltage source: Circuit symbol and

Independent Current Source (ICS)

1.1. CIRCUIT INTRODUCTORY

Figure 1.9: Independent current source: Circuit symbol and

Figure 1.10: Short circuit: Circuit symbol and i-v characteristic.

Figure 1.11: Open circuit: Circuit symbol and i-v characteristic.

CHAPTER 1. FUNDAMENTALS OF CIRCUIT ANALYSIS

Consider the circuit as in Figure 1.12. Find i1

We rst mark the circuit variables consisting of 11 unknowns:

1.1. CIRCUIT INTRODUCTORY

identify the nodes, and write the i-v characteristics equations:

Note that we have 6 elements so we have 12 circuit variables. However,

Note: As can be seen, even relatively small circuits results in a large

CHAPTER 1. FUNDAMENTALS OF CIRCUIT ANALYSIS

Circuit Reduction Techniques

Combination of KVLs, KCLs, and i-v characteristics equations result in

1.2. CIRCUIT REDUCTION TECHNIQUES

A box with some

Figure 1.13: Subcircuit and rest of the circuit.

CHAPTER 1. FUNDAMENTALS OF CIRCUIT ANALYSIS

Subcircuits play an important role in linear circuit theory. Thevenin

Figure 1.14: Two resistors

Substituting from i-v characteristics equations in KVL, we get

1.2. CIRCUIT REDUCTION TECHNIQUES

one resistor with

A Resistor in Series with a Current Source

Figure 1.15: A resistors in

Figure 1.16: Thevenin form.

A Resistor in Series with a Voltage Source

CHAPTER 1. FUNDAMENTALS OF CIRCUIT ANALYSIS

A Voltage Source in Series with a Voltage Source

A Voltage Source in Series with a Current Source

1.2. CIRCUIT REDUCTION TECHNIQUES

KCL : i = is1 = is2