ENCYCLOPAEDIA OF MATHEMATICS

Volume 1
 ENCYCLOPAEDIA OF MATHEMATICS

Managing Editor

M. Hazewinkel

Scientific Board

J. F. Adams, S. Albeverio, J. B. Alblas, S. A. Amitsur, I. J. Bakelman, J. W. de Bakker,

C. Bardos, H. Bart, M. Bercovier, M. Berger, E. A. Bergshoeff, E. Bertin, F. Beukers,
A. Beutelspacher, K.-D. Bierstedt, H. P. Boas, J. Bochnak, H. J. M. Bos, B. L. J. Braaksma,
A. E. Brouwer, M. G. de Bruin, R. G. Burns,Bums, P. Cartier, J. M. C. Clark, Ph. Clement,
A. M. Cohen, J. W. Cohen, R. F. Curtain, M. H. A. Davis, M. V. Dekster, C. Dellacherie,
G. van Dijk, H. C. Doets, I. Dolgachev, A. Dress, L. P. D. van den Dries, J. J. Duistermaat,
D. van Dulst, H. van Duyn, H. Dym, A. Dynin, W. Eckhaus, J. Eells, P. van Emde Boas,
H. Engl, G. Eskin, G. Ewald, V. I. Fabrikant, A. Fasano, M. Fliess, R. M. Fossum,
M. Freidlin, B. Fuchssteiner, G. B. M. van der Geer, R. D. Gill, V. V. Goldberg, J. de Graaf,
J. Grasman, P. A. Griffiths, L. Gross, P. Gruber, E. J. Hannan, K. P. Hart, G. Heckmann,
A. J. Hermans, W. H. Hesselink, C. C. Heyde, K. Hirsch, M. W. Hirsch, A. T. de Hoop,
P. J. van der Houwen, N. M. Hugenholtz, C. B. Huijsmans, A. Isidori, E. M. de Jager,
P. T. Johnstone, D. Jungnickel, M. A. Kaashoek, V. Kac, W. L. J. van der Kallen,
D. Kanevsky, Y. Kannai, H. Kaul, M. S. Keane, E. A. de Kerf, W. Klingenberg, T. Kloek,
C. Koeman, J. A. C. Kolk, G. Komen, T. H. Koornwinder, Koomwinder, S. Kotz, L. Krop,
B. Kuperschmidt, H. A. Lauwerier, J. van Leeuwen, J. Lennox, H. W. Lenstra Jr.,
J. K. Lenstra, H. Lenz, M. Levi, J. Lindenstrauss, J. H. van Lint, A. Liulevicius, M. Livshits,
W. A. J. Luxemburg, R. M. M. Mattheij, L. G. T. Meertens, P. Mekenkamp, A. R. Meyer,
J. van Mill, I. Moerdijk, H. Neunzert, G. Y. Nieuwland, G. J. Olsder, F. van Ostayen,
B. Pareigis, K. R. Parthasarathy, I. I. Piatetskii-Shapiro, H. G. J. Pijls, N. U. Prabhu,
G. B. Preston, E. Primrose, A. Ramm, C. M. Ringel, J. B. T. M. Roerdink, K. W. Roggenkamp,
G. Rozenberg, W. Rudin, F. H. Ruymgaart, S. N. M. Ruysenaars, A. Salam, A. Salomaa,
P. Saunders, C. L. Scheffer, R. Schneider, J. A. Schouten, A. Schrijver, F. Schurer,
I. A. Segal, J. J. Seidel, A. Shenitzer, V. Snaith, T. A. Springer, J. H. M. Steenbrink,
J. D. Stegeman, F. W. Steutel, P. Stevenhagen, I. Stewart, L. Streit, K. Stromberg,
L. G. Suttorp, D. Tabak, F. Takens, R. J. Takens, N. Temme, S. H. Tijs, B. Trakhtenbrot,
N. S. Trudinger, L. N. Vaserstein, M. L. J. van de Vel, F. D. Veldkamp, W. Vervaat,
P. M. B. Vitanyi, H. A. van der Vorst, J. de Vries, F. Waldhausen, B. Wegner,
J. J. O. O. Wiegerinck, J. Wiegold, J. C. Willems, J. M. Wills, B. de Wit, S. A. Wouthuysen,
S. Yuzvinskii, L. Zalcman, S. I. Zukhovitzkii
ENCYCLOPAEDIA
OF
MATHEMATICS
Volume 1
A-B

An updated and annotated translation of the Soviet

'Mathematical Encyclopaedia'

Reidel
KLUWER ACADEMIC PUBLISHERS
Science and Technology Division
Dordrecht / Boston / Lancaster / Tokyo
 Library of Congress Cataloging-in-Publication Data

Matematicheskaia entsiklopediia. English.

Encyclopaedia of mathematics.

1. Mathematics- - Dictionaries. I. Hazewinkel, Michie!. II. Title.

QA5.M3713 1987 510'.3'21 87-26437

Published by Reidel
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 SOVIET MATHEMATICAL ENCYCLOPAEDIA

Editor-in-Chief

I. M. Vinogradov

Editorial Board

S. 1. Adyan, P. S. Aleksandrov, N. S. Bakhvalov, A. V. Bitsadze, V. 1. Bityutskov (Deputy

Editor-in-Chief), L. N. BOl'shev, A. A. Gonchar, N. V. Efimov, V. A. Il'in, A. A. Karatsuba,
L. D. Kudryavtsev, B. M. Levitan, K. K. Mardzhanishvili, E. F. Mishchenko, S. P. Novikov,
E. G. Poznyak, Yu. V. Prokhorov (Deputy Editor-in-Chief), A. 1. Shirshov, A. G. Sveshnikov,
A. N. Tikhonov, P. L. UI'yanov, S. V. Yablonskii

Translation Arrangements Committee

V. I. Bityutskov, R. V. Gamkrelidze, Yu. V. Prokhorov

'Soviet Encyclopaedia' Publishing House

 PREFACE

This ENCYCLOPAEDIA OF MATHEMATICS aims to be a reference work for all parts of mathema-
tics. It is a translation with updates and editorial comments of the Soviet Mathematical Encyclo-
paedia published by 'Soviet Encyclopaedia Publishing House' in five volumes in 1977 - 1985.
The annotated translation consists of ten volumes including a special index volume.
There are three kinds of articles in this ENCYCLOPAEDIA. First of all there are survey-type
articles dealing with the various main directions in mathematics (where a rather fine subdivision
has been used). The main requirement for these articles has been that they should give a reason-
ably complete up-to-date account of the current state of affairs in these areas and that they
should be maximally accessible. On the whole, these articles should be understandable to
mathematics students in their first specialization years, to graduates from other mathematical
areas and, depending on the specific subject, to specialists in other domains of science, en-
gineers and teachers of mathematics. These articles treat their material at a fairly general level
and aim to give an idea of the kind of problems, techniques and concepts involved in the area in
question. They also contain background and motivation rather than precise statements of pre-
cise theorems with detailed definitions and technical details on how to carry out proofs and con-
structions.
The second kind of article, of medium length, contains more detailed concrete problems, re-
sults and techniques. These are aimed at a smaller group of readers and require more back-
ground expertise. Often these articles contain more precise and refined accounts of topics and
results touched on in a general way in the first kind of article.
Finally, there is a third kind of article: short (reference) definitions.
Practically all articles (all except a few of the third kind) contain a list of references by means
of which more details and more material on the topic can be found. Most articles were specially
written for the encyclopaedia and in such cases the names of the original Soviet authors are
mentioned. Some articles have another origin such as the Great Soviet Encyclopaedia (Bol,' shaya
Sovieticheskaya Entsiklopediya or BSE).
Communication between mathematicians in various parts of the world has certainly greatly
improved in the last decennia. However this does not mean that there are so-to-speak 'one-to-
one onto' translations of the terminology, concepts and tools used by one mathematical school
to those of another. There also are varying traditions of which questions are important and
which not, and what is considered a central problem in one tradition may well be besides the
point from the point of view of another. Even for well established areas of mathematical inquiry,
terminology varies across languages and even within a given language domain. Further, a concept,
theorem, algorithm, ... , which is associated with one proper name within one tradition may
well have another one in another; especially if the result or idea in question was indeed disco-
vered independently and more-or-Iess simultaneously. Finally, mathematics is a very dynamic
science and much has happened since the original articles were finalized (mostly around 1977).
This made updates desirable (when needed). All this, as well as providing additional references
to Western literature when needed, meant an enormous amount of work for the board of

vii
PREFACE

experts as a whole; some indeed have done a truly impressive amount of work. I must stress
though that I am totally responsible for what is finally included and what is not of all the mate-
rial provided by the members of the board of experts.
Many articles thus are provided with an editorial comment section in a different and some-
what smaller typeface. In particular, these annotations contain additional material, amplifica-
tions, alternative names, additional references, .... Modifications, updates and other extra
material provided by the original Soviet authors (not a rare occurrence) have been incorporated
in the articles themselves.
The final (lO-th) volume of this ENCYCLOPAEDIA OF MATHEMATICS will be an index volume.
This index will contain all the titles of the articles (some 6600) and in addition the names of all
the definitions, named theorems, algorithms, lemmas, scholia, constructions, ... , which occur
in the various articles. This includes, but is by no means limited to, all items which are printed
bold or italic. Bold words or phrases, by the way, always refer to another article with (precisely)
that title.
All articles have been provided with one or more AMS classification numbers according to
the 1980 classification scheme (not, for various reasons, the 1985 revision). So have all items
occurring in the index. A phrase or word out of an article which is included in the index always
inherits all the classification numbers of the article in question. In addition, it may have been
provided with its own classification numbers. In the index volume these numbers will be listed
with the phrase in question. Thus e.g. the Quillen - Suslin theorem of algebraic K-theory will
have its own main classification numbers (these are printed in bold; in this case that number is
18F25) as well as a number of others, often from totally different fields, pointing e.g. to parts of
mathematics where the theorem is applied, or where there occurs a problem related to it. In this
case e.g. 93D15. The index volume will also contain the inversion of this list which thus for each
number will provide a list of words and phrases which may serve as an initial description of the
'content' of that classification number (as far as this ENCYCLOPAEDIA is concerned). For more
details on the index volume, its structure and organisation, and what kind of things can be done
with it, d. the (future) special preface to that volume.
Classifying articles is a subjective matter. Opinions vary greatly as to what belongs where and
thus this attempt will certainly reflect the tastes and opinions of those who did the classification
work. One feature of the present classification attempt is that the general basic concepts and de-
finitions of an area like e.g. 55N (Homology and Cohomology theories) or 601 (Markov proces-
ses) have been assigned classification numbers like 55NXX and 60JXX if there was no finer clas-
sification number different from ... 99 to which it clearly completely belongs.
Different parts of mathematics tend to have differences in notation. As a rule, in this
ENCYCLOPAEDIA in a given article a notation is used which is traditional in the corresponding
field. Thus for example the (repeated index) summation convention is used in articles about
topics in fields where that is traditional (such as in certain parts of differential geometry (tensor
geometry)) and it is not used in other articles (e.g. on summation of series). This pertains espe-
cially to the more technical articles.
For proper names in Cyrillic the British Standards Institute transcription system has been
used (cf. Mathematical Reviews). This makes well known names like S. N. Bernstein come out
like Bernshte'in.
In such cases, especially in names of theorems and article titles, the traditional spelling has
been retained and the standard transcription version is given between brackets.
Ideally an encyclopaedia should be complete up to a certain more-or-Iess well defined level of
detail. In the present case I would like to aim at the completeness level that every theorem,

viii
 PREFACE

concept, definition, lemma, construction which has a more-or-Iess constant and accepted name
by which it is referred to by a recognizable group of mathematicians occurs somewhere, and can be
found via the index. It is unlikely that this completeness ideal will be reached with this present
ENCYCLOPAEDIA OF MATHEMATICS, but it certainly takes substantial steps in this direction.
Everyone who uses this ENCYCLOPAEDIA and find items which are not covered, which, he feels,
should have been included, is invited to inform me about it. When enough material has come in
in this way supplementary volumes will be put together.
The ENCYCLOPAEDIA is alphabetical. Many titles consist of several words. Thus the problem
arises how to order them. There are several systematic ways of doing this of course. For inst-
ance using the first noun. All are unsatisfactory in one way or another. Here an attempt has
been made to order things according to words or natural groups of words as they are daily used
in practice. Some sample titles may serve to illustrate this: Statistical mechanics, mathematical
problems in; Lie algebra; Free algebra; Associative algebra; Absolute continuity; Abstract
algebraic geometry; Boolean functions, normal forms of. Here again taste plays a role (and us-
ages vary). The index will contain all permutations. Meanwhile it will be advisable for the read-
er to tryout an occasional transposition himself. Titles like K-theory are to be found under K,
more precisely its lexicographic place is identical with 'K theory'. I.e. '-' = 'space' and comes
before all other symbols. Greek letters come before the corresponding Latin ones, using the
standard transcriptions. Thus X2-distribution (chi-squared distribution) is at the beginning of the
letter C. A * as in C* -algebra and *-regular ring is ignored lexicographically. Some titles involve
Greek letters spelled out in Latin. These are of course ordered just like any other 'ordinary'
title.
This volume has been computer typeset using the (Unix-based) system of the CWI, Am-
sterdam. The technical (mark-up-Ianguage) keyboarding was done by Rosemary Daniels,
Chahrzade van 't Hoff and Joke Pesch. To meet the data-base and typesetting requirements of this
ENCYCLOPAEDIA substantial amounts of additional programming had to be done. This was done by
lohan Wolleswinkel. Checking the translations against the original texts, and a lot of desk edit-
ing and daily coordination was in the hands of Rob Hoksbergen. All these persons, the mem-
bers of the board of experts, and numerous others who provided information, remarks and
material for the editorial comments, I thank most cordially for their past and continuing efforts.
The original Soviet version had a printrun of 150.000 and is completely sold out. I hope that
this annotated and updated translation will turn out to be comparably useful.

Bussum, August 1987 MICHIEL HAZEWINKEL

ix
 A
A -INTEGRAL - One of the generalizations of the countable intersection, and is idempotent. With respect
Lebesgue integral, given by E. Titchmarsh [1] for the to ~ -operations, the Haire property (of subsets of an
integration of functions conjugate to summable ones. A arbitrary topological space) and the property of being
measurable function f(x) is called A -integrable over Lebesgue measurable are invariant.
[a, b] if
m{x: If(x) 1>n} = 0 !
[ ) References
[IJ Al.EKsANDRov, P.S.: c.R. Acad. Sci. Paris 162 (1916), 323-325.
[2J Al.EKsANDROV, P.S.: Theory offunctions of a real variable and
and if the theory of topological spaces, Moscow, 1978 (in Russian).
b [3] KOI.MOGOROV, A.N.: 'P.S. Aleksandrov and the theory of as-
I = lim j[f(x)ln fix operations', Uspekhi Mat. Nauk 21, no. 4 (1966), 275-278 (in
n~oo a
Russian).
exists, where [4] SusLIN, M.YA.: c.R. Acad. Sci. Paris 164 (1917),88-91.
[5] LUZIN, N.N.: Selected works, Vol 2, Moscow, 1958, p. 284.
f(X) if 1 f(x) 1 "';;n, [6J KURATOWSICI, K..: Topology, Acad. Press, 1966-1968 (translated
[f(x)ln = {0 if 1 f(x) 1 >n. from the French). A. G. El'kin
The number I is called the A -integral. It is denoted by Editorial comments. The d -operation is in the West
b usually attributed to M.Ya. Suslin [4], and is therefore also
(A) jf(x)dx. called the Suslin operation, the Suslin d -operation or the
d
Suslin operation d. d -sets are usually called analytic sets.
References AMS 1980 Subject Classification: 04A15
[IJ 1iTcHMARSH, E.G.: 'On conjugate functions', Proc. London
Math. Soc. 29 (1928),49-80.
[2J VINOGRADOVA, IA.: 'Generalized integrals and Fourier series', A POSTERIORI DISTRIBUTION - A conditional
ftogi Nauk. Mat. Anal. 1970 (1971),65-107 (in Russian).
probability distribution of a random variable, to be
LA. Vinogradova contrasted with its unconditional or a priori distribu-
AMS 1980 Subject Classification: 26A42 tion.
Let 8 be a random parameter with an a priori den-
.s;/-oPERATION, operation ~ - A set-theoretical sity p(8), let X be a random result of observations and
operation, discovered by P.S. Aleksandrov [1] (see also let p(x 18) be the conditional density of X when 8=0;
[2], [3D. Let {E", ... lit} be a system of sets indexed by then the a posteriori distribution of 8 for a given
all finite sequences of natural numbers. The set X = x, according to the Bayes formula, has the density
co
P = U nEil, ...... ' p(fJ 1x) = +}(8)p(x 18)
" .. ___ ..... __ k=J
j p{fI)p(x l8)dfJ
where the union is over all infinite sequences of natural -co
numbers, is called the result of the ~ -operation If T(x) is a sufficient statistic for the family of distribu-
applied to the system {E", ... lit}. tions with densities p (x I 8), then the a posteriori distri-
The use of the ~-operation for the system of inter- bution depends not on x itself, but on T(x). The
vals of the number line gives sets (called ~ -sets in asymptotic behaviour of the a posteriori distribution
honour of Aleksandrov) which need not be Borel sets p(Olx), ... ,x,,) as n~oo, where Xj are the results of
(see ~ -set; Descriptive set theory). The ~ -operation is independent observations with density p(x I(0 ), is
stronger than the operation of countable union and 'almost independent' of the a priori distribution of 8.

1
A POSTERIORI DISTRIBUTION
For the role played by a posteriori distributions in
the theory of statistical decisions, see Bayesian
approach.
References
[1] BERNSHTEiN, S.N.: Probability theory, Moscow-Leningrad, 1946
(in Russian). v V.. P rokh orov
~u.
Editorial comments.
References
[A1] SVERDRUP, E.: Laws and chance variations, 1, North-Holland,
1967, p. 214ft.
AMS 1980 Subject Classification: 62F15
A POSTERIORI PROBABILITY of an event - The
conditional probability of an event taking place under
certain conditions, to be contrasted with its uncondi-
tional or a priori probability. There is no difference
between the meaning of the terms 'conditional' and 'a
posteriori'. The former term is employed if the condi-
tion itself is hypothetical and is not directly observed in
the course of the experiment. The latter term is
employed if it is desired to stress that the condition in
question is actually observed. The a posteriori probabil-
ity is connected with the a priori probability by the
Bayes fonnula.
Yu. V. Prokhorov
AMS 1980 Subject Classification: 62F15
A PRIORI DISTRIBUTION - The probability distri-
bution of a random variable, to be contrasted with the
conditional distribution of this random variable under
certain additional conditions. Usually the term 'a priori
distribution' is used in the following way. Let (8, X) be
a pair of random variables (random vectors or more
general random elements). The random variable 8 is
considered to be unknown, while X is considered to be
the result of an observation to be used for estimation of
e. The joint distribution of e and X is given by the
distribution of 8 (now called the a priori distribution)
and the set of conditional probabilities Po of the ran-
dom variable X given e=o. According to the Bayes
fonnula, one can calculate the conditional probability
of 8 with respect to X (which is now called the a pos-
teriori distribution of 8). In statistical problems, the a
priori distribution is often unknown (and even the
assumption on its existence is not sufficiently founded).
For the use of the a priori distribution, see Bayesian
approach.
Yu. V. Prokhorov
Editorial comments.
References
[A1] SVERDRUP, E.: Laws and chance variations, 1, North-Holland,
1967, p. 214ft.
AMS 1980 Subject Classification: 62F15
2
A PRIORI PROBABILITY of an event - The proba-
bility of the event, considered in contrast to the condi-
tional probability of the same event under certain addi-
tional conditions. The latter is called the a posteriori
probability. This terminology is, as a rule, used in con-
nection with the Bayes fonnula.
Yu. V. Prokhorov
AMS 1980 Subject Classification: 60A 15
sI-SET, analytic set, in a complete separable metric
space - A continuous image of a Borel set. Since any
Borel set is a continuous image of the set of irrational
numbers, an d -set can be defined as a continuous
image of the set of irrational numbers. A countable
intersection and a countable union of d -sets is an d-
set. Any d -set is Lebesgue-measurable. The property
of being an d -set is invariant relative to Borel-
measurable mappings, and to d -operations (cf. d -
operation). Moreover, for a set to be an d -set it is
necessary and sufficient that it can be represented as
the result of an d -operation applied to a family of
closed sets. There are examples of d -sets which are
not Borel sets; thus, in the space 21 of all closed sub-
sets of the unit interval I of the real numbers, the set of
all closed uncountable sets is an d -set, but is not
Borel. Any uncountable d -set topologically contains a
perfect Cantor set. Thus, d -sets 'realize' the contin-
uum hypothesis: their cardinality is either finite, No or
2No. The Luzin separability principles hold for d -sets.
References
[1] KURATOWSKl, K.: Top%gy, 1, Acad. Press, 1966-1968
(translated from the French).
[2] LUZIN, N.N.: Lectures on analytic sets and their applications,
Moscow, 1953 (in Russian).
B.A. Efimov
Editorial comments. Nowadays the class of analytic sets
is denoted by };1, and the class of co-analytic sets (cf.
Cd-set) by 111.
References
[A1] JECH, T.J.: The axiom of chOice, North-Holland, 1973.
[A2] MOSCHOVAKIS, Y.N.: Descriptive set theory, North-Holland,
1980.
AMS 1980 Subject Classification: 04A 15, 54H05
A -SYSTEM - A 'countably-ramified' system of sets,
i.e. a family {Ani'" n,} of subsets of a set X, indexed
by all finite sequences of natural numbers. An A-
system {Ani" . nJ IS called regular if
Ani' .. n,n CAn n,'
I ... A sequence
HI
Ani' .. ,Ani "n,'" of elements of an A-system
indexed by all segments of one and the same finite
sequence of natural numbers is called a chain of this
A-system. The intersection of all elements of a chain is
called its kernel, and the union of all kernels of all

chains of an A -system is called the kernel of this A-
system, or the result of the .s;/ -operation applied to this
A -system, or the .s;/ -set generated by this A -system.
Every A-system can be regularized without changing
the kernel (it suffices to put
A~, "'n, =AI n '" nAn, n'). If At is a system of
sets, then the kernels of the A -system composed from
the elements of At are called the .s;/-sets generated by
At . The d-sets generated by the closed sets of a topo-
logical space are called the .s;/-sets of this space.
References
[1] ALEKSANDROV, P.S.: An introduction to set theory and general
topology, Moscow, 1977 (in Russian).
[2] KURATOWSKI, K.: Topology, 1, Acad. Press, 1966 (translated
from the French).
A.G. El'kin
Editorial comments. Thed-operation is an important
tool in descriptive set theory. It was introduced by M.Ya.
Suslin, hence it is also known as the Sus/in operation (also
So us/in operation). In this connection, an A-system is also
called a Sus/in (Souslin) scheme. See also d-operation;
d-set.
While [2] is the standard reference for classical results, a
modern approach can be found in [A1].
References
[A 1] CHRISTENSEN, J.P.R.: Topology and Borel structure, North-
Holland, 1974.
AMS 1980 Subject Classification: 54H05, 04A 15
ABACUS - 1) A calculating frame used for arithmeti-
cal calculations in ancient Greece, ancient Rome and
later, up to the 18-th century, in Western Europe. The
frame was subdivided into strips, inside which the
counters (bones, pebbles, etc.) were moved. The
Chinese equivalent of the abacus - swan-pan - was
used in the Far East, while the 'schety' was employed
in Russia.
2) In nomography an abacus is a special sketch (the
so-called net chart nomogram).
BSE-3
Editorial comments.
References
[A 1] ROUSE BALL, W. W.: A short account of the history of
mathematics, Dover, reprint, 1960, pp. 123-125.
AMS 1980 Subject Classification: OOA08, OOA25,
01A20
ABEL CRITERION - 1) Abel's criterion for series of
numbers. If the series
n=l
is convergent and if the numbers an form a monotone
bounded sequence, then the series
ABEL DIFFERENTIAL EQUATION
is convergent.
2) Abel's criterion for series offunctions. The series
n=l
converges uniformly on a set X if the series
n=1
converges uniformly on X and if the functions an(x),
n = 1,2, ... , for any x EX, form a monotone sequence
that is uniformly bounded on X An Abel criterion for
the uniform convergence of integrals
00
ja(n, x)b(n, x)dn,
which depend on a parameter x EX, can be formulated
in a similar manner.
The Abel criteria can be strengthened (see, for exam-
ple, Dedekind criterion (convergence of series. See
also Dirichlet criterion (convergence of series); Abel
transformation.
References
[1] FICHTENHOLZ, G.M.: Differential und Integralrechnung, I,
Deutsch. Verlag Wissenschaft., 1964.
[2] KUDRYAVTSEV, L.D.: Mathematical analysis, I, Moscow, 1973
(in Russian).
[3] WHITTAKER, E.T. and WATSON, G.N.: A course of modern
analysis, 1-2, Cambridge Univ. Press, 1952.
L.P. Kuptsov
AMS 1980 Subject Classification: 40A05, 40A10
ABEL DIFFERENTIAL EQUATION - The ordinary
differential equation
y' = !O(X)+!I(X)Y +/2(x)y 2 +!3 (x)y 3
(Abel's differential equation of the first kind) or
[go(x)+gl(X)Y}y' = !O(X)+!I(X)y+/2(x)y 2+/J(X)y 3
(Abel's differential equation of the second kind). These
equations arose in the context of the studies of N.H.
Abel [1] on the theory of elliptic functions. Abel's dif-
ferential equations of the first kind represent a natural
generalization of the Riccati equation.
If fl EC(a, b), fz and h ECI(a, b) and h(x)i=O for
x E[a, b], then Abel's differential equation of the first
kind can be reduced to the normal form
dz / dt=z3 +cI>(t) by substitution of variables [2]. In
the general case, Abel's differential equation of the first
kind cannot be integrated in closed form, though this is
possible in special cases [2]. If gO,gIECI(a,b) and
gl(X)i=O, gO(X)+gl(X)yi=O, Abel's differential equa-
tion of the second kind can be reduced to Abel's dif-
ferential equation of the first kind by substituting
go(x)+gl(x)y= 1 / z.
Abel's differential equations of the first and second
kinds, as well as their further generalizations
3
ABEL DIFFERENTIAL EQUATION

n m n Let A~ be the class of functions 1 (z) of the form

y' = ~(x>'y;, y'~gj(x>'yj = ~(x>.y;, co _
;=0 j=O ;=0
f(z) = ~(n!)-"a"zn, lim I a" Illn :E;; T,
have been studied in detail in the complex domain (see, n=O n-+oo

for example, [3]).
References
[I] ABEL, N.H.: 'Precis d'une th60rie des fonctions elliptiques', J.
Reine Angew. Math. 4 (1829), 309-348.
[2] K.uoo!, E.: Differentialgleichungen: LOsungsmethoden und
O<r < 00, O~a< 00, and let Aa. be the class of all pos-
sible sequences {)..,,} such that I)"" I ~(n + 1 I , r-
n =0, 1, .... By definition, the bound 01 convergence for
the class Aa. is the supremum S of the values of r for
/I(

LOsungen, 1, Chelsea, reprint, 1971. which every function 1 (z) EA~ can be represented by a
[3] GoLUBEV, V.V.: Vorlesungen iiber Differentialgleichungen im series (*). The infimum W of the r-values for which
/I(

Komplexen, Deutsch. Verlag Wissenscbaft., 1958 (translated there exist a function I(Z)EA~ and a sequence
from the Russian). N.Kh. Rozov {)..,,}EA.. such thatj<n)()..,,)=O, n=O, 1, ... , l(z)"O,
is called the bound 01 uniqueness. The quantities WI
AMS 1980 Subject Classification: 34AXX
and Wo are known, respectively, as the Whittaker con-
stant and the Goncharov constant. It has been shown
ABEL - GoNCHAROV PROBLEM, Goncharov prob- that S I = WI (cf. [6]); the more general statements
lem - A problem in the theory of functions of a com-
plex variable, consisting of finding the set of all func- S .. = WI> W.. = WI> O:E;;a<oo,
tions 1 (z) from some class satisfying the relations have also been proved [5], [10].
j<n)()..,,)=A n (n=O, 1, ... ). Here {An} and {)..,,} are Thus, if {~n}EA.., the Abel-Goncharov problem is
sequences of complex numbers admissible for the given reduced to finding the constant WI. Its precise numeri-
class. The problem was posed by V.L. Goncharov [2]. cal value is not known, but one has obtained the
A function 1 (z) is put in correspondence with the bounds 0.7259< WI <0.7378 [9].
series Relative to the Abel- Goncharov problem for the
""
~j<n)(An)Pn(z), (*) class A'i of functions which are regular in the region
n=O
I z I ;;;;.1 and such that I( 00) = 0, the following has been
the Abel-Goncharov interpolation series, where Pn(z) is shown: For any set of numbers {~}, I ~ I;;;;'1, which
the Goncharov polynomial defined by the equalities satisfy the condition
P~k)(At) = 0, k=O, ... ,n-I; p~n)(z) - l. . nk
N.H. Abel gave a formal treatment of the case k~1A;f = 0,
)""=a+nh, n=O, 1, ... , where a and h are real where {nk} is an ~creas~~ ~ence ::.f natural
numbers (h=j=O) (cf. [1]). Here numbers, the equatIons f
(~k)-O, k-O, 1, ... ,
Pn(z) = -\-(z-a)(z-a-nhr- I imply l(z)=O. Moreover, for any number b>O there
n. exist a sequence {~n}
The series (*) is used to study the zeros of the succes-
Jim_n_ =b
sive derivatives of regular functions. The set of func- I An
n .... "" I
tions 1 (z) representable by a series (*) is called the con- and a function l(z) .. O,/(z)EAi, for which
vergence class of the Abel- Goncharov problem.
J<n)(~n)=O, n =0, 1, . .. [7].
In the case when limn . . . oo I ~n I = 00, the The Abel- Goncharov problem includes the so-called
Abel- Goncharov convergence class has been expressed two-point problem posed by J.M. Whittaker [12]. Let the
in terms of bounds on the order and the type of entire sequences {Jlk } and {p.,. } be such that
functions I(z), depending on the growth of the quan- n
{Jld U {p.,.} = {n} and {Jld {p.,.} = 0. The problem
tity ~~:~ I ~k+ I - ~k I [2]. is to determine conditions under which there exists a
If ).." =nl/Pl(n), where l(n) is a slowly increasing function l(z)"O, re~ar on the segment [0, 1], satisfy-
function, O<p< 00, the Abel- Goncharov convergence ing the conditions f r.)(1) =0, r)(O) =0. This problem
class has, in a sense, been exactly determined [6]. Other was solved for various subclasses of the class of func-
Abel- Goncharov convergence classes for entire func- tions regular in the disc I z I<R (R> 1). The resulting
tions of finite and infinite order have been identified in conditions, which are in a sense precise, are expressed
terms of various constraints on the indicators of the in terms of various bounds imposed on the coefficients
respective classes of functions. The Abel- Goncharov a,. of the expansions
problem has also been studied for entire functions of
00
several variables. Exact estimates of Goncharov polyno- j(z) = ~a z
mials were obtained for some classes of interpolation k=O

nodes. depending on {Jld [3]. This problem was generalized

4
 ABEL INTEGRAL EQUATION

and solved using methods of the theory of infinite sys- Abel's inequality is proved by means of the Abel
tems of linear equations [4]. In the particular case transfonnation.
where the sequence {pd forms an arithmetic progres- LD. Kudryavtsev
sion and one is dealing with entire functions of AMS 1980 Subject Classification: 26015
exponential type, the two-point problem has, in a cer-
tain sense, been solved completely [8]. ABEL INTEGRAL EQUATION - The integral equa-
References tion
[1] ABEL, N.H.: 'Sur les fonctions generatrices et leurs
determinants', in Oeuvres completes, Vol. 2, Christiania, 1839,
I
x cp(s) ds = f(x),
Yx-s
(1)

pp. 77-88. Edition de Holmboe.
[2] GoNTCHAROFF, W.L. [V.L. GoNCHAROV]: 'Recherches sur les
derivees successives des fonctions analytiques. Generalisation
de la serie d'Abel', Ann. Sci. Ecole Norm. Sup. (3) 47 (1930),
1-78.
[3] GEL'FOND, A.O. and IBRAGIMOV, I.I.: 'On functions, the
derivatives of which vanish in two points', Izv. Akad. Nauk
SSSR Ser. Mat. 11 (1947), 547-560 (in Russian).
[4] DZHRBASHYAN, M.M.: 'Uniqueness and representability
theorems for analytic functions', Izv. Akad. Nauk SSSR Ser.
Mat. 16 (1952), 225-252 (in Russian).
[5] DRAGILEV, M.M. and CHuIrnLOVA, O.P.: 'On convergence of
certain interpolation series', Sibirsk. Mat. Zh. 4, no. 2 (1963),
287-294 (in Russian).
[6] EVGRAFov, M.A.: The Abel-Goncharov interpolation problem,
to which the solution of the Abel problem is reduced.
The generalized Abel integral e~tion is the equation
Ja(x-sf
cp(s) ds = f(x), a';;;x ~b, (2)
where a>O and O<a<l are known constants, f(x) is
a known function and <p(x) is the unknown function.
The expression (x - s) - a is called the kernel of the
Abel integral equation, or Abel kernel. An Abel
integral equation belongs to the class of Volterra equa-
tions of the first kind (d. Volterra equation). The equa-
tion b

falx-sla
Moscow, 1954 (in Russian).
[7] KAz'MIN, YU.A.: 'On zeros of sequences of derivatives of ana- p(s) ds = f(x), a~x~b, (3)
lytic functions', Vestnik Moskov. Univ. Ser. I Math. Mekh., no.
1 (1963), 26-34 (in Russian). is called Abel's integral equation with fixed limits.
[8] KAz'MIN, Yu.A.: 'On a problem of Gel'fond- Ibragimov',
If f (x) is a continuously-differentiable function, then
Vestnik Moskov. Univ. SeT. I Math. Mekh. 19, no. 6 (1965),
37-44 (in Russian). the Abel integral equation (2) has a unique continuous
[9] MAcINTYRE, S.S.: 'On the zeros of successive derivatives of solution given by the formula
integral functions', Trans. Amer. Math. Soc. 67 (1949),241-251.
[10] SUETIN, Yu.K.: 'Coincidence of constants of uniqueness and cp(x) = sina'IT ~fx (t)dt (4)
convergence of certain interpolation problems', Vestnik 'IT dx a (x-t)l-a'
Moskov. Univ. Ser. I Math. Mekh. 21, no. 5 (1966), 16-25 (in
Russian). or
[11] OsKOLKOV, V.A.: 'On estimates for Gonwov polynomials', cp(x) = sina'IT [ (a) +fx t)dt ] (5)
Math. USSR-Sb. 21, no. 1 (1973), 57-62. (Mat. Sb. 92, no. 1 'IT (x-a)l-a a (x-t)l-a '
(1973), 55-59)
[I2] WHITTAKER, J.M.: Interpolatory function theory, Cambridge which is the same thing. Formula (5) also gives the
Univ. Press, 1935. solution of an Abel integral equation (2) under more
V.A. Oskolkov
general assumptions [3], [4]. Thus, it has been shown [3]
Editorial comments. For an introduction to the general that if f (x) is absolutely continuous on the interval
area, cf. [A 1]. [a, b], the Abel integral equation (2) has the unique
References solution in the class of Lebesgue-integrable functions
[A1] BoAS, RP.: Entire functions, Acad. Press, 1954. given by formula (5). For the solution of the Abel
AMS 1980 Subject Classification: 30E05 integral equation (3) see [2]; see also [6].
References
ABEL INEQUALITY - An estimate for the sum of [1] BOCHER, M.: 'On the regions of convergence of power-series
which represent two-dimensional harmonic functions', Trans.
products of two numbers. If sets of numbers ak and bk
Amer. Math. Soc. 10 (1909),271-278.
are given such that the absolute values of all sums [2] CARLEMAN, T.: 'Ueber die Abelsche Integralgleichung mit kon-
Bk=b) + ... +bko k=l, ... ,n, are bounded by a stanten Integrationsgrenzen', Math. Z. 15 (1922),111-120.
number B, i.e. I Bk I ~B, and if either ai ;;:'ai +) or [3] ToNELU, L.: 'Su un problema di Abel', Math. Ann. 99 (1928),
183-199.
ai ~ai + J, i = 1, ... ,n - 1, then [4] TAMAllKIN, J.D.: 'On integrable solutions of Abel's integral

Ik~lakbk I.;;; B(I al 1+21 an I).

equation', Ann. of Math. (2) 31 (1930), 219-229.
[5] MnrnuN, S.G.: Linear integral equatiOns, Hindushtan Publ.
Comp., Delhi, 1960 (translated from the Russian).
If the ak are non-increasing and non-negative, one [6] GAKHOV, F.D.: Boundary value problems, Pergamon, 1%6
has the simpler estimate: (translated from the Russian). B V. Kh -1 l'd
.. veue I ze

Ik~lakbk I.;;; Bal' Editorial comments. The left-hand side of (2) is also
known as a Riemann - Liouville fractional integral, where

5
ABEL INTEGRAL EQUATION

Rea<1, cf. [A1]. If one integrates in (1) and (2) from x to
00 and replaces x - s by s - x, then one obtains left-hand
sides which are known as Abel transform and Weyl frac-
tional integral (ct. [A1)), respectively. This Abel transform
is the case SL(2, R) of the Abel transform on a real semi-
simple Lie group, cf. [A2].
A general source for integral equations is [A3].
References
[A1] ERDELYI, A., MAGNUS, W., OBERHETINGER, F. and TRIcOMl,
F.G.: Tables of integral transforms, II, McGraw-Hili, 1954,
Chapt. 13.
[A2] GoDEMENT, R.: 'Introduction aux travaux de A. Selberg', in
sem. Bourbaki, Vol. 144, 1957.
[A3] FENYO, S. and STOLLE, H.W.: Theorie und Praxis der linearen
Integralgleichungen, 3, Birkhauser, 1984, Sect. 13.2.4.
AMS 1980 Subject Classification: 45E10
ABEL - POISSON SUMMATION MEfHOD - One
of the methods for summing Fourier series. The Fourier
series of a function fEL[O,271"] is summable by the
Abel- Poisson method at a point cp to a number S if
lim j(p, </ = S,
as M(p, </ approaches P(1, </>0) non-tangentially inside the
disc, cf. [A2], pp. 129-130.
References
[A1] ZYGMUND, A.: Trigonometric series, 1-2, Cambridge Univ,
Press, 1979.
[A2] Tsun, M.: Potential theory in modern function theory,
Maruzen, 1959.
AMS 1980 Subject Classification: 40G 10
ABEL PROBLEM - To find, in a vertical plane (s, r),
a curve such that a material point moving along it
under gravity from rest, starting from a point with
ordinate x, will meet the r-axis after a time T = f (x),
where the function f (x) is given in advance. The prob-
lem was posed by N.H. Abel in 1823, and its solution
involves one of the first integral equations - the Abel
integral equation - which was also solved. In fact, if w
is the angle formed by the tangent of the curve being
sought with the r-axis, then
~ =- V2g(x -s) sinw.

where
p-->i-O
Integrating this equation between and x and putting
aO _.1_ = <t>(s), - -y2g<I>(x) = j(x),
= 2+ ~(akcosk</>+bksink</l,
00
j(p,</ smw
k=i
one obtains the integral equation
ji( "') =
p, 'I'
1- j" ji("'+ t)
'1T _" 'I'
1- p2 dt
2(1-2pcost+p2)'
(*)
f p(s)ds
o Vx-s
= j(x)
If fEC(O, 271"), then the integral on the right-hand side
for the unknown function cp(s), the determination of
is a harmonic function for I z I I pe i </> I <1, which is,
as has been shown by S. Poisson, a solution of the Diri- which makes it possible to find the equation of the
chlet problem for the disc. The Abel summation method curve being sought. The solution of the equation intro-
applied to Fourier series was therefore named the duced above is:
Abel- Poisson summation method, and the integral (*)
was named the Poisson integral.
<t>(x) =~
'1T
[1(Ql+ f ((T)dr].
Vx VX-T
0

If (p, cp) are polar coordinates of a point inside the

disc of radius one, then one can consider the limit of
f(p, cp) as the point M(p, cp) approaches a point on the
bounding circle not by a radial or by a tangential but
rather along an arbitrary path. In this situation the
Schwarz theorem applies: If fbelongs to L[O, 271"] and is
continuous at a point </>0, then
(P,</>~,</>o/(p, </
irrespective of the path along which the point M(p, cp)
approaches the point P(l, </>0) as long as that path
remains within the disc with radius one.
References
[I] BARY, N.K. [N.K. BARI]: A treatise on trigonometric series, Per-
gamon, 1964 (translated from the Russian).
Editorial comments. A theorem related to Schwarz'
theorem stated above is Fatou's theorem: If fEL[O, 2'1T], then
for al most all </>0
lim
(p, </-->(1, </>0)
f(p, </
References
[1] ABEL, N.H.: 'Solutions de quelques problemes Ii l'aide
d'integrales defines', in Oeuvres completes, nouvelle Iid., Vol. I,
Grondahl and Son, Christiania, 1881, pp. 11-27.
B. V Khvedelidze
Editorial comments. In the case that f(x)=const, this is
the famous tautochrone problem first solved by Chr,
Huyghens, who showed that this curve is then a cycloid.
= j(</>o)
References
[A 1] JERRI, A.J.: Introduction to integral equations with applica-
tions, M. Dekker, 1985, Sect. 2.3,
[A2] HOCHSTADT, H.: Integral equations, Wiley (Interscience),
1973.
[A3] MOISEIWITSCH, B.L.: Integral equations, Longman, 1977.
A A Z kh AMS 1980 Subject Classification: 70D05, 45E10
. . a arov
ABEL SUMMATION METHOD - One of the
methods for the summation of series of numbers. The
series
= f(</>o)

6

can be summed by the Abel method (A -method) to the
number S if, for any real x, O<x < 1, the series
is convergent and
This summation method can already be found in the
works of L. Euler and G. Leibnitz. The name 'Abel
summation method' originates from the Abel theorem
on the continuity of power series. The Abel summation
method belongs to the class of totally regular summa-
tion methods and is more powerful than the entire set
of Cesaro summation methods. The Abel summation
method is used in conjunction with Tauberian theorems
to demonstrate the convergence of a series.
A closely related summation method is the A-
method. Let z be a complex number, I z I < 1; the
series 00
~ak
k=O
is summed by the A -method to the number S if
00
lim ~akzk = S,
k=O
where z~ 1 along any path not tangent to the unit cir-
cle.
References
[I] HARDy, G.H.: Divergent series, Clarendon, 1949.
[2] BARY, N.K. [N.K. BARI]: A treatise on trigonometric series, Per-
gamon, 1964 (translated from the Russian). A A Z kh
. . a arov
AMS 1980 Subject Classification: 40G10
ABEL THEOREM - 1) Abel's theorem on algebraiC
equations: Formulas expressing the solution of an arbi-
trary equation of degree n in terms of its coefficients
using radicals do not exist for any n ;:'5. The theorem
was proved by N.H. Abel in 1824 [1]. Abel's theorem
may also be obtained as a corollary of Galois theory,
from which a more general theorem follows: For any
n;:' 5 there exist algebraic equations with integer coeffi-
cients whose roots cannot be expressed in terms of rad-
icals of rational numbers. For a modern formulation of
Abel's theorem for equations over an arbitrary field, see
Algebraic equation.
2) Abel's theorem on power series: If the power series
00
S(z) = ~ak(z-bt, (*)
k=O
where ab b, z are complex numbers, converges at
z = z 0, then it converges absolutely and uniformly
within any disc I z -b I ~p of radius p< I Zo -b I and
with centre at b. The theorem was established by N.H.
Abel [2]. It follows from the theorem that there exists a
ABEL THEOREM
number R E[O, 00] such that if I z - b I <R the series is
convergent, while if I z - b I > R the series is divergent.
The number R is called the radius of convergence of the
series (*), while the disc I z - b I <R is known as the
disc of convergence of the series (*).
3) Abel's continuity theorem: If the power series (*)
converges at a point z 0 on the boundary of the disc of
convergence, then it is a continuous function in any
closed triangle T with vertices z 0, Z 1 , z 2, where z J, z 2
are located inside the disc of convergence. In particu-
lar
lim S(z) = S(ZO).
z-+zo,
ZET
This limit always exists along the radius: The series (*)
converges uniformly along any radius of the disc of
convergence joining the points b and z o. This theorem
is used, in particular, to calculate the sum of a power
series which converges at the boundary points of the
disc of convergence.
4) Abel's theorem on Dirichlet series: If the Dirichlet
series
.p(S) = ~ane-'\"s, s=a+it, An>O,
n=1
converges at the point So =00 +ito, then it converges in
the half-plane 0>00 and converges uniformly inside
any angle I arg(s-so) I ~O<7T /2. It is a generaliza-
tion of Abel's theorem on power series (take An =n and
put e -s = z). It follows from the theorem that the
domain of convergence of a Dirichlet series is some
half-plane o>c, where c is the abscissa of convergence
of the series.
The following theorem is valid for an ordinary Diri-
chlet series (when An =lnn) with a known asymptotic
behaviour for the sum-function An =al + ... +an of
the coefficients of the series: If
An = B nSI(lnn)a +O(n P),
where B, s 1 ,0: are complex numbers, (3 is a real
number, 01-1<{3<01, 01 =Resl, then the Dirichlet
series converges for 01 <0, and the function cp(s) can be
regularly extended to the half-plane (3<0 with the
exception of the point s = s I. Moreover
.p(s) = B f(u+ l)s(s -sd- a - I +g(s)
if 0:*-1, -2, ... ,and
AJ) -- B (_u_l)!ss-si
'i'\S (_l}-a (
ns-si )+ gs
)-a-Il ( ()
if 0: = - 1, - 2, . . .. Here g (s) is a regular function for
0>{3.
E.g., the Riemann zeta-function t{s) (An =n, B = 1,
s 1 = I, 0: =0, (3)0) is regular at least in the half-plane
0>0, with the exception of the point s = 1 at which it
has a first-order pole with residue 1. This theorem can
be generalized in various ways. E.g., if
k
An = ~B)nSj(lnntj +O(n P),
)=1
7
ABEL THEOREM
where Bj , Sj' aj (l:S;;;;j:S;;;;k), are arbitrary complex
numbers, and Ok -1 <{3<Ok < ... <01, then the Diri-
chlet series converges for 0>0J, and <Pes) is regular in
the domain 0> {3 with the exception of the points
S 1, . ,Sk at which it has algebraic or logarithmic
singularities. Theorems of this type yield information
on the behaviour of the Dirichlet series in a given half-
plane, based on the asymptotic behaviour of An.
References
[1] ABEL, N.H.: Oeuvres completes, Vol. 1, Christiania, 1881.
[2] ABEL, N.H.: 'Untersuchungen tiber die Reihe
l+mx /2+m(m-l)x 2 /(2,1)+", ',J. Reine Angew. Math.
I (1826), 311-339.
[3] MARKUSHEVICH, A.!.: Theory offunctions of a complex variable,
1, Chelsea, 1977 (translated from the Russian).
L.P. Kuptsov
Editorial comments. More on Abel's theorems 2) - 4)
can be found in [A 1].
References
[A1] HARDY, G.H.: Divergent series, Clarendon, 1949.
AMS 1980 Subject Classification: 10HXX, 12AXX,
12E12, 30830, 30815,40G10
ABEL TRANSFORMATION, summation by parts - A
transformation
N N-'
~akbk = aNBN-a,B o - ~Bk(ak+,-ak)'
k=' k='
where ak> bk are given, Bo is arbitrarily selected, and
Bk = Bk-,+bk =Bo+b,+ ... +bb k=I, ... ,N.
The Abel transformation is the discrete analogue of the
formula for integration by parts.
If aW"~O and if the sequence {Bk} is bounded, then
the Abel transformation can be applied to the series
00 00
~akbk = ~(ak -ak+,)Bk -a,Bo
k=' k='
The Abel transformation is used to prove several cri-
teria of convergence of series of numbers and functions
(cf. Abel criterion). The Abel transformation of a series
often yields a series with an identical sum, but with a
better convergence. It is also regularly used to obtain
certain estimates (cf. Abel inequality), in particular, for
investigations on the rate of convergence of a series. It
was introduced by N.H. Abel [1].
References
[I] ABEL, N.H.: 'Untersuchungen tiber die Reihe
I+mx /2+m(m-1)x 2 /(2'1)+'" ',l. Reine Angew. Math.
1(1826),311-339.
[2] MARKUSHEVICH, A.I.: Theory offunctions of a complex variable,
I, Chelsea, reprint, 1977 (translated from the Russian).
[3] WHITTAKER, E.T. and WATSON, G.N.: A course of modern
analysis, 1-2. Cambridge Univ. Press, 1952.
L.P. Kuptsov
AMS 1980 Subject Classification: 65810
ABELIAN CATEGORY - A category that displays
8
some of the characteristic properties of the category of
all Abelian groups. Abelian categories were introduced
as the basis for an abstract construction of homological
algebra [4]. A category ~ is said to be Abelian [2] if it
satisfies the following axioms:
AO. A null object exists (cf. Null object of a
category).
AI. Each morphism has a kernel (d. Kernel of a
morphism in a category) and a cokernel.
A2. Each monomorphism is a normal monomorphism,
so that it occurs as the kernel of a morphism; each epi-
morpliism is a normal epimorphism.
A3. A product and a coproduct exist for each pair of
objects (cf. Product of a family of objects in a category).
In defining an Abelian category it is often assumed
that ~ is a locally small category (d. Small category).
The coproduct of two objects A and B of an Abelian
category is also know as the direct sum of these objects
and is denoted by A EB B, A liB or A -+ B.
Examples of Abelian categories. 1) The dual category
of an Abelian category is also an Abelian category.
2) The category R we of all unitary left modules over
an arbitrary associative ring R with a unit element and
all R-module homomorphisms is an Abelian category
(e.g. the category of all Abelian groups).
3) Any full subcategory of an Abelian category, that
is, a subcategory that contains for each one of its mor-
phisms also the kernel and cokernel of that morphism,
and that contains for each pair of objects A and B also
their product and coproduct, is an Abelian category.
The small Abelian categories are exhausted by the
subcategories of the above type of categories R we of
unitary left modules. Accordingly, the following
Mitchell theorem is valid: For each small Abelian
category there exists a full exact imbedding into some
category R we.
4) Any category of diagrams i5{'tI, 21), with diagram
scheme 'tI over an Abelian category ~, is an Abelian
category. In the scheme 'V one may distinguish the set
C of commutativity relations, i.e. the set of pairs (cp, !f;)
of paths CP=(CP1,'" ,CPn), !f;=(!f;I,'" ,!f;m) in 'tI with a
common begin and end. Then the complete subcategory
of the category i5{'tI, ~) generated by all those diagrams
D : 'tI~~ that satisfy
D(</ = D(</>,) ... D(</>n) = D(l/;,) ... D(l/;m) = D(l/;)
is an Abelian category. In particular, if 1) is a small
category and if the set C consists of all pairs of the
form (a/3, y) with y=a/3, then the corresponding sub-
category is the Abelian category of one-place covariant
functors from 'tI into ~ (cf. Functor).
Suppose that a null object exists in a small category
1). A functor F: '1)~~ will then be called normalized if
it takes a null object into a null object. The complete

subcategory of the category of functors generated by
the normalized functors will then be an Abelian
category. In particular, if ~ is a category the objects of
which are all integers and the null object is N, while the
non-null non-identity morphisms form a sequence
L, do d,
"'~(-1)~O~1~2~'"
in which dndn + 1 =0, n=O, -+-1, -+-2, ... , then the
corresponding subcategory generated by the normalized
functors is called the category of complexes over ~. On
the category of complexes there are defined the additive
functors zn, B n , H n of n-dimensional cycles, n-
dimensional boundaries and n-dimensional homologies,
respectively, with values in~. These constitute the
basis for the development of homological algebra.
5) A complete subcategory ~I of an Abelian category
~ is said to be dense if it contains all subobjects (d.
Subobject) and quotient objects (d. Quotient object) of
its objects, and if for the exact sequence
O~A~B---,>C~O
BEOb~1 if and only if A, CEOb~I' The quotient
category ~ / ~I is then constructed as follows. Let
(R, J.t] be a subobject of the direct sum A EB B with pro-
jections '171, '172 and let the square
R ---'> A
/L'172 ~ ~a
B ---'> X categories and functors, Wiley, 1968.
fJ
be co-universal (i.e. a co-fibred product). The subobject
(R, J.t] is called an ~I-subobject if Cokeq.L'17I, Kerj3E~I'
Two ~I-subobjects are equivalent if they contain some
~I-subobject. By definition, the set H iJI / iJI , (A, B) con-
sists of the equivalence classes of ~I-subobjects. The
usual multiplication of binary relations is compatible
with the equivalence relation thus introduced, which
makes it possible to construct the quotient category
~ / ~I' which is an Abelian category. The faithful
functor T: ~~~ / ~I is defined by assigning to each
morphism ex:A-'?B its graph in A EBB. A subcategory
~I is said to be a localizing subcategory if T has a com-
plete univalent right-adjoint functor Q: ~ / ~I -'?~.
6) For any topological space X the category of left
G-modules over X, where G is a sheaf of rings with unit
element over X, is an Abelian category.
It is possible to introduce into any Abelian category
~ a partial sum of morphisms so that A becomes an
additive category. For this reason the product and the
coproduct of any pair of objects in an Abelian category
are identical. Moreover, in defining an Abelian
category it suffices to assume the existence of either
products or coproducts. Any Abelian category is a
bicategory with a unique bicategorical structure. These
properties characterize an Abelian category: A category
ABELIAN DIFFERE]\iTIAL
~ with finite products is Abelian if and only if it is
additive and if any morphism ex has a kernel and a co-
kernel and can be decomposed into a product
a = Coker(Kera) 8 ker(Cokera),
in which () is an isomorphism.
The Mitchell theorem quoted above constitutes the
underlying principle of the so-called 'diagram-chasing'
method in an Abelian category: Any proposition about
commutative diagrams that is valid for all categories of
left modules R fIR and that is a consequence of the
exactness of certain sequences of morphisms, is valid in
all Abelian categories.
In a locally small Abelian category, the ~-subobjects
of an arbitrary object form a Dedekind lattice. If pro-
ducts (or coproducts) of any family of objects exist in
~, this lattice will be complete. These conditions are
known to be met if there is a generating object U in ~
and if the coproducts
II
iEJ
Vj, V, = V,
exist for any set I. These conditions are satisfied, for
instance, by Grothendieck categories (d. Grothendieck
category), which are equivalent to the quotient
categories of the category of modules by a localizing
subcategory (the Gabriel- Popescu theorem).
References
[1] BUCUR, I. and DELEANU, A.: Introduction to the theory of
[2] FREYD, P.: Abelian categories: An introduction to the theory of
functors, Harper and Row, 1964.
[3] GABRIEL, P.: 'Des categories Abeliennes', Bull. Soc. Math.
France 90 (1962), 323-448.
[4] GROTHENDIECK, A.: 'Sur quelques points d'algebre homo1o-
gique', Tohoku Math. J. 9 (1957), 119-221. M Sh T. I k
. . sa en 0
Editorial comments. Composition of morphisms is writ-
ten from left to right in this article; i.e. # denotes the com-
position of <jl: A~B, >J;: B~C. A dense subcategory is more
often called a Serre subcategory.
References
[A 1] MITCHELL, 8.: Theory of categories, Acad. Press, 1965.
[A2] POPESCU, N.: Abelian categories with applications to rings
and modules, Acad. Press, 1973.
AMS 1980 Subject Classification: 18EXX
ABELIAN DIFFERENTIAL - A holomorphic or
meromorphic differential on a compact, or closed,
Riemann surface S (d. Differential on a Riemann sur-
face).
Let g be the genus of the surface S (d. Genus of a
surface); let alb I ... agbg be the cycles of a canonical
basis of the homology of S. Depending on the nature
of their singular points, one distinguishes three kinds of
Abelian differentials: I, II and III, with proper inclu-
sions I clI CIII. Abelian differentials of the first kind
are first-order differentials that are holomorphic every-
9
ABELIAN DIFFERENTIAL

where on S and that, in a neighbourhood U of each a_ n

[- + ... a_I]
+-+f(z) dz, (3)
point POES, have the form w=pdz=p(z)dz, where zn z
z = x + iy is a local uniforrnizing variable in U, where f (z) is a regular function, n is the order of the
dz =dx +i dy, and p(z) is a holomorphic, or regular,
pole (if a -n=l=O), and a - I is the residue of the pole. If
analytic function of z in U. The addition of Abelian n = 1, the pole is said to be simple. An Abelian dif-
differentials and multiplication by a holomorphic func- ferential of the second kind is a meromorphic differen-
tion are defined by natural rules: If tial all residues of which are zero, i.e. a meromorphic
w = pdz, 'TT = qdz, a = a(z), differential with local representation
then
w+'TT = (p+q)dz, aw = (ap)dz. [ a_
n + ... + a-2 +f(Z)]dZ,
zn Z2
The Abelian differentials of the first kind form a g-
An Abelian differential of the third kind is an arbitrary
dimensional vector space 21. After the introduction of
the scalar product Abelian differential.
Let w be an arbitrary Abelian differential with A-
(w, 'TT) = f f w*"iT,
periods A I, . . . ,Ag; the Abelian differential
s
w' = w - A I </>1 - ... - Ag</>g then has zero A -periods
where w*7i is the exterior product of w with the star-
and is known as a normalized Abelian differential. In
conjugate differential 7i, the space 21 becomes a Hilbert
particular, if PI and P 2 are any two points on S, one
space.
can construct a normalized Abelian differential wI,2
Let A I, B I, . . . ,Ag, Bg be the A- and B-periods of
the Abelian differential of the first kind w, i.e. the with the singularities (1 / z)dz in PI and (-1/ z)dz in
P 2 , which is known as a normal Abelian differential of
integrals
the third kind. Let w be an arbitrary Abelian differential
Aj = fw, Bj = fw, j=l, ... ,g. with residues c I, . . . ,Cn at the respective points
aJ hJ
PI,'" ,Pn ; then, always, CI + ... +cn=O. If Po is
The following relation then holds:
any arbitrary point on S such that Po=l=Pj ,
II w 112 = ((A)ij-BjA;) ;;;. O. (1) ) = 1, ... ,n, then w can be represented as a linear
j=\ combination of a normalized Abelian differential of the
If A'I, B'I, ... ,A~, B~ are the periods of another second kind W2, a finite number of normal Abelian dif-
Abelian differential of the first kind 'TT, then one has ferentials of the third kind Wj,O, and basis Abelian dif-
ferentials of the first kind </>k:
i(w, "iT) = (A;B) - BjA)) = O. (2)
j=\
w = W2 + c;w;,o + Ak<l>k'
The relations (1) and (2) are known as the bilinear j = \ k =\

Riemann relations for Abelian differentials' of the first Let W3 be an Abelian differential of the third kind with
kind. A canonical basis of the Abelian differentials of only simple poles with residues Cj at the points Pj'
the first kind, i.e. a canonical basis </>1, ... , </>g of the ) = 1, ... , n, and let WI be an arbitrary Abelian dif-
space 21, can be chosen so that ferential of the first kind:
Aij = f<l>l = B,; , Ak = fW\, Bk = fW\,
Q) Qk hk

where 8;; = 1 and 8ij =0 if )=I=i. The matrix (B;), A~ = fW3' B~ = fW3' k=l, ... ,g,
i,) = 1, ... ,g, of the B-periods Ok hA,

where the cycles ak, bk do not pass through the poles of

BIj = f<l>l
b, W3' Let the point Po ES not lie on the cycles ab bk and
IS then symmetric, and the matrix of the imaginary let L j be a path from Po to Pj' One then obtains bi-
parts (1m B;j) is positive definite. An Abelian differen- linear relations for Abelian differentials of the first and
tial of the first kind for which all the A-periods or all third kinds:
the B-periods are zero is identically equal to zero. If all
the periods of an Abelian differential of the first kind w (AkB~-BkA~)
=I
= 27Td>JWl'
are real, then w = O. k ; - I l'j

Abelian differentials of the second and third kinds are,
in general, meromorphic differentials, i.e. analytic dif-
ferentials which have on S not more than a finite set of
singular points that are poles and which have local
representations
Bilinear relations of a similar type also exist between
Abelian differentials of the first and second kinds.
In addition to the A- and B-periods Ab Bb
k = I, ... ,g, known as the cyclic periods, an arbitrary
Abelian differential of the third kind also has polar

10

periods of the form 27Ticj along zero-homologous cycles
which encircle the poles Pj' One thus has, for an arbi-
trary cycle 'I,
fW3 = (lkA k+ lg+k B k)+27Tim/j,
y k=1 ;=1
where h, 19+k' and mj are integers.
Important properties of Abelian differentials are
described in terms of divisors. Let (w) be the divisor of
the Abelian differential w, i.e. (w) is an expression of
the type (w) = p~l ... P~', where the Prs are all the
zeros and poles of wand where the ars are their multi-
plicities or orders. The degree d[(w)]=a] + ... +an of
the divisor of the Abelian differential w depends only
on the genus of S, and one always has d[(w)]=2g-2.
Let a be some given divisor. Let n(a) denote the com-
plex vector space of Abelian differentials w of which
the divisors (w) are multiples of a, and let L(a) denote
the vector space of meromorphic functions f on S of
which the divisors (f) are multiples of a. Then
dimQ(a)=dimL(a / (w)). Other important information
on the dimension of these spaces is contained in the
Riemann - Roch theorem: The equality
dimL(a -I )-dim~(a) = d[a]- g + 1
is valid for any divisor a. It follows from the above, for
example, that if g = 1, i.e. on the surface of a torus, a
meromorphic function cannot have a single simple pole.
Let S be an arbitrary compact Riemann surface on
which there are meromorphic functions z and w which
satisfy an irreducible algebraic equation F(z, w) = O.
Any arbitrary Abelian differential on S can then be
expressed as w=R(z, w)dz where R(z, w) is some
rational function in z and w; conversely, the expression
w=R(z, w)dz is an Abelian differential. This means
that an arbitrary Abelian integral
fR(z, w)dz = fw
is the integral of some Abelian differential on a com-
pact Riemann surface S.
See also Algebraic function.
References
[I] SPRINGER, G.: Introduction to Riemann surfaces, Addison-
Wesley, 1957, Chapt.lO.
[2] NEVANLINNA, R.: Uniformisierung, Springer, 1953. Chapt.5.
[3] CHEBOTAREV, N.G.: The theory of algebraic functions,
Moscow-Leningrad, 1948, Chapt.3;8 (in Russian).
E.D. Solomentsev
Editorial comments. Another good reference, replacing
[3], is [A1].
References
[A1] LANG, S.: Introduction to algebraic and abelian functions,
Addison-Wesley, 1972.
AMS 1980 Subject Classification: 14H05, 30F30
ABELIAN FUNCflON - A generalization of the con-
ABELIAN FUNCTION
cept of an eUiptic function of one complex variable to
the case of several complex variables. A function f (z)
in the variables z 1, . . . ,Zp' Z =(z 1, . . . ,zp), which is
meromorphic in the complex space a, p;;;. I, is called
an Abelian function if there exist 2p row vectors in a
which are linearly independent over the field of real
numbers and are such that fez +w.)= fez) for all
ZEa, p=I, ... ,2p. The vectors w. are called the
periods or the system of periods of the Abelian function
f (z). All periods of the Abelian function f (z) form an
Abelian group r under addition, which is known as the
period group (or the period module). A basis of this
group is known as a basis system of periods of the
Abelian function, or also as a system of basic (or primi-
tive) periods. An Abelian function fez) is said to be
degenerate if there exists a linear transformation of the
variables z], ... ,Zp which converts f (z) into a func-
tion of fewer variables; otherwise f (z) is said to be a
non-degenerate Abelian function. Degenerate Abelian
functions are distinguished by having infinitely small
periods, i.e. for any number >0 it is possible to find a
period w. for which
II wp II = Vk~1 I wkp 12 < {.
If p = I, the non-degenerate Abelian functions are ellip-
tic functions of one complex variable. Each Abelian
function with period group r is naturally identified
with a meromorphic function on the complex torus
a/ r, i.e. the quotient space a/
r (cf. Quasi-Abelian
function).
The study of Abelian functions began in the 19-th
century in connection with the inversion of Abelian
integrals (cf. Abelian integral) of the first kind (cf.
Jacobi inversion problem, [1], [2]). The Abelian func-
tions obtained in solving this problem are called speCial
Abelian functions; sometimes, in earlier work, only such
functions were considered to be Abelian functions. Let
u], ... ,up be linearly independent normal Abelian
integrals of the first kind, constructed on a Riemann
surface F:
x x
= f du I, . . , up = f dup'
U1
c] Cp
and let
XI Xp
Zj = fdu;+ ... + fdu;, )=1, ... ,p,
Cl Cp
be a given system of sums in which the lower integra-
tion limits c 1, . . . ,cp are considered as fixed on the
surface F. It is then possible to define the special
Abelian functions as all rational functions in the p
coordinates of the upper limits x], ... ,xp ' the latter
11
 ABELIAN FUNCTION
being in turn considered as functions of p points
z 1, . . . ,zp on F. In the symbolic notation introduced
by K. Weierstrass, any special Abelian function Al(z)
may be represented as
Al(z) = Al(zb ... ,zp) =
= R[X\(Zb ... ,zp), ... ,XP(Zb ... ,zp)].
The complex tori a / r corresponding to special
Abelian functions are the Jacobi varieties of algebraic
curves.
The matrix W whose columns form a period basis of
the Abelian function f (z) has dimension p X 2p and is
known as the period matrix of the Abelian function
f (z). A necessary and sufficient condition for a given
matrix W of dimension p X 2p to be the period matrix
of some non-degenerate Abelian function f (z) is for
the matrix to satisfy the following conditions (the
Riemann- Frobenius conditions). It must be a Riemann
matrix, i.e. there must exist an anti-symmetric non-
degenerate square matrix M with integer elements, of
order 2p, and such that 1) WMW T =0, where WT is
the transposed matrix of W; and 2) the matrix
iWMW T defines a positive-definite Hermitian form [3].
If the conditions 1) and 2) are expressed as equations
and inequalities respectively, a system of p (p - 1) / 2
Riemann equations and p (p - 1) / 2 Riemann inequalities
is obtained. The number p is called the genus of the
matrix W and of the corresponding Abelian function
f(z). The columns w v = Re w v + i 1m w v of W, regarded
as vectors in the real Euclidean space R2p, define the
period parallelotope of f (z).
All Abelian functions corresponding to the same
period matrix W form an Abelian function field Kw. If
the field Kw contains a non-degenerate Abelian func-
tion, its degree of transcendence over the field C is p;
the torus 0' / r is then an Abelian variety, and Kw
turns out to be its field of rational functions. If, on the
other hand, all Abelian functions of Kw are degenerate,
then K w is isomorphic to the field of rational functions
on an Abelian variety of dimension lower than p. See
also Quasi-Abelian function.
As in the case of elliptic functions, any Abelian func-
tion can be represented as the quotient of two entire
transcendental theta-functions (d. Theta-function),
which in turn can be represented as theta-series. A
given Riemann matrix W determines a class of theta-
series by means of which all Abelian functions of the
field Kw can be constructed.
For special Abelian functions, the matrix W can
always be brought - by means of a linear transforma-
tion of the independent variables z 1, . . . , Zp - to the
form
12
wi
W=
o
The Riemann relations between the elements of the
matrix "ajk ,I, j, k = 1, ... ,p, now ensure the sym-
metry of the matrix, ajk =akj' and the positive definite-
ness of the matrix of real parts II Re ajk II,
j, k = 1, ... ,p. However, for p>3 there will be only
3(p - 1) independent elements among the elements ajk
of II ajk II, i.e. as many as the number of conformal
moduli for the Riemann surface F on which the inver-
sion problem is solved (cf. Moduli of a Riemann sur-
face). In addition to the Riemann relations there are in
such a case (p - 2)(P - 3) / 2 transcendental relations
between the ajb the explicit form of which for the case
p =4 was first found by E. Schottky in 1886; for a
review of subsequent advances in the field see [5].
Special Abelian functions can be represented as the
quotient of two entire theta-functions with half-integral
characteristics of a special kind. These representations
yield a number of properties of special Abelian func-
tions which generalize many properties of elliptic func-
tions. Thus, derivatives of an Abelian function f (z)
with respect to any argument Zj are Abelian functions;
any p + 1 Abelian functions are related by some alge-
braic equation; any Abelian function can be rationally
expressed in terms of some p + 1 Abelian functions, e.g.
by means of an arbitrary Abelian function and its p
partial derivatives of the first order; Abelian functions
satisfy the addition theorem, i.e. the value of an Abelian
function at the point a + b E a may be rationally
expressed in terms of the values of a certain (p + 1)-
tuple of Abelian functions at the points a, b EO'.
Abelian functions are very important in algebraic
geometry as a means of uniformization of algebraic
varieties of certain types.
References
[1] SPRINGER, G.: Introduction to Riemann surfaces, Addison-
Wesley, 1957.
[2] CHEBOTAREV, N.G.: The theory of algebraic functions,
Moscow-Leningrad, 1948 (in Russian).
[3] SIEGEL, c.L.: Automorphe Funktionen in mehrerer Variablen,
Math. Inst. Gbttingen, 1955.
[4] WElL, A.: Introduction a l'hude des varihes kahleriennes, Her-
mann, 1958.
[5] MUMFORD, D.: Curves and their Jacobians, Univ. of Michigan
Press, 1975.
[6] STAHL, H.: Theorie der Abelschen Funktionen, Leipzig, 1896.
[7] CLEBSCH, A. and GORDAN, P.: Theorie der Abelschen Funk-
tionen, Teubner, 1866.
[8] CONFORTO, F.: Abelsche Funktionen und algebraische
Geometrie, Springer, 1956.
E.D. Solomentsev
Editorial comments. A basis system of periods is also
known as a period basis. Instead of [5] one can also con-

suit the more recent [A 1]. The Riemann - Frobenius condi-
tions are also known as the Riemann bilinear relations or
simply as the Riemann relations. The Schottky problem con-
cerns the description ot the space of all Jacobians in the
space of matrices satisfying the Riemann conditions. It has
recently been solved, ct. Jacobi variety.
References
[A1] ARBELLO, E., CORNALBA, M., GRIFFITHS, P.A. and HARRIs, 1.:
Geometry of algebraic curves, 1, Springer, 1985.
AMS 1980 Subject Classification: 33A25, 32A20
ABELIAN GROUP - A group whose operation is
commutative (cf. Commutativity). They are named after
N.H. Abel, who used such groups in the theory of solv-
ing algebraic equations by means of radicals. It is cus-
tomary to write the operation in an Abelian group in
additive notation, i.e. to use the plus sign (+ ) for that
operation, called 'addition', and the zero sign (0) for
the neutral element, named zero (in multiplicative nota-
tion it is called the unit element).
Examples of Abelian groups. All cyclic groups (cf.
Cyclic group) are Abelian, in particular, the additive
group of integers. All direct sums (cf. Direct sum) of
cyclic groups are Abelian. Also the additive group of
rational numbers Q is Abelian; it is moreover a locally
cyclic group, i.e. a group in which all finitely generated
subgroups are cyclic. Finally, the groups of type p 00
(or the quasi-cyclic groups Zp~ ), where p is an arbi-
trary prime number, are Abelian (cf. Group of type poo;
Quasi-cyclic group).
The free composition in the variety of Abelian groups
coincides with the direct sum. A free Abelian group is a
direct sum of infinite cyclic groups. Every subgroup of
a free Abelian group is free Abelian. The set of all ele-
ments of finite order in an Abelian group forms a sub-
group, which is called the torsion subgroup (periodic
part) of the Abelian group. The quotient group of an
Abelian group by its torsion subgroup is a group
without torsion. Thus, every Abelian group is an exten-
sion of a torsion-free group by a torsion group. This
extension does not always split, i.e. the torsion group is
usually not a direct summand. An Abelian torsion
group in which the order of every element is a power of
a fixed prime number p is said to be primary with
respect to p (the term p-group is used in general group
theory). Every torsion group splits uniquely into a
direct sum of primary groups that correspond to dis-
tinct prime numbers.
A complete classification is known for finitely gen-
erated Abelian groups. This is given by the fundamental
theorem of finitely generated Abelian groups: Every fin-
itely generated Abelian group is a direct sum of finitely
many non-split cyclic subgroups some of which are fin-
ABELIAN GROUP
ite and primary, while the others are infinite (G. Fro-
benius, L. Stickelberger). In particular, finite Abelian
groups split into a direct sum of primary cyclic groups.
Such splittings are, in general, not unique, but any two
splittings of a finitely generated Abelian group into
direct sums of non-split cyclic groups are isomorphic,
so that the number of infinite cyclic summands and the
collection of the orders of the primary cyclic summands
do not depend on the splittings chosen. These numbers
are called invariants of the finitely generated Abelian
group. They constitute a complete system of invariants,
in the sense that two (finitely generated) Abelian
groups are isomorphic if and only if they have the same
invariants. Each subgroup of a finitely generated
Abelian group is itself finitely generated.
Not every Abelian group is a direct sum of (even
infinitely many) cyclic groups. For primary groups
there is a necessary and sufficient condition for the
existence of such splittings, the so-called Kulikov cri-
terion. Let A be a primary Abelian group for a prime
number p. A non-zero element a of A is said to be an
element of infinite height in A if for any integer k the
equation p k X = a is solvable in A, and a is called an ele-
ment of height n if this equation is solvable for k ~n
only. The Kulikov citerion states that a primary Abelian
group A is a direct sum of cyclic groups if and only if
it is the union of an increasing sequence of subgroups
such that the set of heights of the elements in each one
of these subgroups is bounded. Any subgroup of an
Abelian group that is a direct sum of cyclic groups
itself is such a direct sum. The indecomposable (non-
splittable) into direct sums primary Abelian groups are
the primary cyclic groups and the group Zp~.
A finite set of elements gl, ... , gk in an Abelian
group is called linearly dependent if there exist integers
n I, . . . , nk, not all equal to zero, such that
~~=lnigi=O. If such numbers do not exist, the set is
said to be linearly independent. An arbitrary collection
of elements of A is said to be linearly dependent if
there exists a finite linearly dependent subcollection.
An Abelian group that is not a torsion group has maxi-
mal linearly independent sets. The cardinality of all
maximal linearly independent collections of elements is
the same and is called the rank (Profer rank) of the
given Abelian group. The rank of a torsion group is
considered to be zero. The rank of a free Abelian group
coincides with the cardinality of a set of free generators
of it.
Every torsion-free Abelian group of rank 1 is iso-
morphic to some subgroup of the additive group of
rational numbers. There exists a complete description
of such groups in the language of types. Each element
of an Abelian group without torsion determines a
13
ABELIAN GROUP
characteristic, which is a countable sequence consisting
of non-negative numbers and the symbol 00. This
sequence is constructed in the following manner. Sup-
pose that the prime numbers are enumerated in increas-
ing order PI.P2, . ... An element a determines the
sequence whose k-th element is equal to Sk if the equa-
tion a =I'i: x is solvable in the group but the equation
a = pi: + x is not solvable, and the k-th element is equal
to 00 if a =PkX is solvable for all s. Two characteristics
are equivalent if they are equal except possibly for a
fInite number of places and if the symbol 00 occurs in
exactly the same places in both sequences. The charac-
teristics of two linearly dependent elements are equal.
A class of equivalent characteristics is said to be a type.
To each torsion-free Abelian group of rank I there
corresponds a uniquely determined type, called the type
of the given group; non-isomorphic groups having dif-
ferent types.
A torsion-free Abelian group that splits into a direct
sum of groups of rank I is said to be completely split.
Not every subgroup of a completely split group is
completely split, though every direct summand is. For
each integer n there exists a torsion-free Abelian group
of rank n that cannot be decomposed (or split) into a
direct sum. For countable torsion-free Abelian groups
one can to construct a complete system of invariants.
An Abelian group is called complete, or divisible, if
for anyone of its elements a and for any integer m the
equation mx =a has a solution in the group. All divisi-
ble Abelian groups turn out to be direct sums of groups
isomorphic to Q and the groups Zp~ , and the cardinal-
ities of the sets of components isomorphic to Q, as well
as to Zp~ (for each p), from a complete and indepen-
dent system of invariants of the divisible group. Every
Abelian groups can be isomorphically imbedded in
some divisible Abelian group. The divisible Abelian
groups and only they are the injective objects in the
category of Abelian groups. A divisible Abelian group
is a direct summand of each Abelian group containing
it. Therefore, an Abelian group is a direct sum of a
divisible Abelian group and a so-called reduced group,
i.e. a group that contains no non-trivial divisible sub-
groups. A classification of reduced Abelian groups is
known only in certain special cases. Thus, VIm's
theorem ([ I]) gives the classification of all countable
reduced Abelian torsion groups.
The theory of Abelian groups has its origins in
number theory, and is now extensively applied in many
modern mathematical theories. Thus, the duality theory
of characters for finite Abelian groups has been consid-
erably extended to the duality theory of locally com-
pact Abelian groups. The development of homological
algebra has made it possible to solve a whole series of
problems in Abelian groups, such as classifying the set
14
of all extensions of one group by another. The theory
of modules is closely connected with Abelian groups
regarded as modules over the ring of integers. Many
results in the theory of Abelian groups can be applied
to the case of modules over a principal ideal ring.
Owing to their relative simplicity and to the fact that
they have been very thoroughly studied (which is con-
firmed, for instance, by the solvability of the elemen-
tary theory of Abelian groups), and to the availability
of a sufficient variety of objects, Abelian groups serve
as a constant source of examples in various fields of
mathematics.
References
[I] KUROSH, AG.: The theory of groups, 1-2, Chelsea, reprint,
1955-1956 (translated from the Russian).
[2] FUCHS, L.: Abelian groups, Hungarian Acad. Sci., Budapest,
1958.
[3] FUCHS, L.: Infinite abelian groups, I, Academic Press, 1970.
[4] KAPLANSKY, I.: Infinite Abelian groups, Michigan Univ. Press,
1954.
[5] Itogi Nauk. Algebra Topol. Geom 1965 (1967), 9-44.
[6] MiSHINA, AP.: 'Abelian groups', J. Soviet Math. 2, no. 3
(1974), 239-263. (Itogi Nauk. i Tekhn. Algebra Topol. Geom 10
(1972),5-45) v
~u.
L E hrs ov
Editorial comments. For an introduction to the theory of
finite rank torsion-free Abelian groups as it has developed
since the publication of the second volume [A2] of [3] in
1973 cf. [A1]. The proceedings [A3] - [A5] give a good idea
of recent results and current research in Abelian groups.
References
[A1] ARNOLD, D.M.: Finite rank torsion free abelian groups and
rings, Lect. Notes in Math., 931, Springer, 1982.
[A2] FUCHS, L.: Infinite abelian groups, 2, Acad. Press, 1973.
[A3] GOBEL, R. and WALKER, E. (EDS.): Abelian group theory,
Lect. Notes in Math., 874, Springer, 1981.
[A4] GOBEL, R., LADY, L. and MADER, A (EDS.): Abelian group
theory, Lect. Notes in Math., 1006, Springer, 1983.
[A5] GOBEL, R., METELLI, c., ORSAITI, A and SALCE, L. (EDS.):
Abelian groups and modulus, CISM courses and lectures,
287, Springer, 1984.
AMS 1980 Subject Classification: 20KXX
ABELIAN INTEGRAL, algebraic integral - An
integral of an algebraic function, i.e. an integral of the
form z]
jR(z, w)dz, (1)
Zo
where R(z, w) is some rational function in variables
z, w that are related by an algebraic equation
F(z, w) - ao(z)w n +al(z)w n -I + ... +an(z) = 0 (2)
with coefficients a/z) that are polynomials in z,
j = 0, ... , n. To equation (2) there corresponds a com-
pact Riemann surface F which is an n-sheeted covering
of the Riemann sphere. On this Riemann surface z, w,
and consequently also R(z, w), can be considered as
single-valued functions of the points on F.
The integral (1) is then given as the integral w of i

the Abelian differential w=R(z, w)dz on F taken along
some rectifiable path L. In general, specifying only the
initial and end points z 0 and z 1 of the path L does not
completely determine the value of the Abelian integral
(1), or, which is the same, the integral (1) turns out to
be a multi-valued function of the initial and end points
of the path L.
The behaviour of an Abelian integral on F depends,
first of all, on the topological structure of F, in particu-
lar on the topological invariant called the genus g of
the surface F (cf. Genus of a surface). The genus g is
connected with the number of sheets n and with the
number of branch points II (counted with multiplicities)
by the relation g = (II / 2) - n + 1. For g = 0 the vari-
ables z and w can be rationally expressed in some
parameter t, and the Abelian integral becomes the
integral of a rational function in t. This will happen,
for example, in the elementary cases w 2 =aOz+al and
w 2 =aoz 2 +alz +a2'
If g:? 1, any Abelian integral can be expressed in the
form of a linear combination of elementary functions
and canonical Abelian integrals of the three kinds. The
integral i w is called an Abelian integral of the first
kind if w is an Abelian differential of the first kind. In
other words, Abelian integrals of the first kind are
characterized by the fact that for a fixed initial point z 0
of the path L they are a function of the upper bound
z 1 that is an everywhere finite, usually multi-valued,
function on F. Such a characterization may be used,
for example, to construct analogues of Abelian integrals
of the first kind on non-compact Riemann surfaces.
Any Abelian integral of the first kind can be
represented in the form of a linear combination of g
linearly-independent normal Abelian integrals of the first
kind
of differentials cf>1, ... ,cf>g that constitute a canonical
basis for the Abelian differentials of the first kind. If
the surface F is cut along the cycles alb 1 . . . agbg of a
canonical basis for the homology, a simply-connected
domain F* is obtained. The integrals i- cf>i are single-
valued functions of the upper bound z 1 for all paths
L * CF* with fixed initial point Zo and fixed end point
z I. The multi-valuedness of the integrals Ui = i cf>i along
an arbitrary path L C F joining z 0 with z 1 is now com-
pletely described by the fact that it differs from the
integrals i- cf>i only by an integral linear combination
of the A-periods Aij and the B-periods Bij of a basis of
the differentials of the first kind. These constitute the
period matrix of dimension g X 2g, which satisfies the
bilinear Riemann relations (cf. Abelian differential).
iL
An integral w where w is an Abelian differential of
ABELIAN INTEGRAL
the second kind is said to be an Abelian integral of the
second kind. Considered as a function of the upper
bound, it has no singularities anywhere on F except for
poles. An Abelian integral of a normalized Abelian
differential of the second kind is known as a normal
Abelian integral of the second kind.
An Abelian integral of the third kind is an arbitrary
Abelian integral. It usually has logarithmic singularities
on F; however, such singularities can only occur in
pairs. An Abelian integral of a normal Abelian dif-
ferential of the third kind is called a normal Abelian
integral of the third kind. Any Abelian integral can be
represented as a linear combination of normal Abelian
integrals of the first, second and third kinds. Unlike
Abelian integrals of the first and second kinds, Abelian
integrals of the third kind usually also have the so-
called polar periods, beside the A- and B-periods (which
are called cyclic periods). Polar periods are taken along
cycles which are homologous to zero, but encircle the
logarithmic singularities of the Abelian integral. They
are caused by the poles of the Abelian differential w
with non-zero residues.
A number of relations depending on the topological
and conformal structure of F exist for arbitrary Abelian
integrals on the same Riemann surface F. Thus, if wP,P J
is a normal Abelian differential of the third kind with
simple poles in Pi and Pi' then the following theorem
on the permutation of the parameters and the bounds of
an Abelian integral of the third kind holds for all
points PI, P 2, P 3, P 4 :
P2 P4
f
PI
WP J P 4 == f
PJ
WP1P2'
Relations which connect Abelian integrals with
rational functions on F are known as Abelian theorems.
In terms of divisors, for example, the Abelian theorem
for Abelian integrals of the first kind has the following
form: A divisor a on F is the divisor of a meromorphic
function if and only if there exists a chain L with
i
aL = a and w = 0 for all Abelian differentials of the
first kind on F. There also exist variants of Abelian
theorems for Abelian integrals of the second and third
kinds [4]. Abelian integrals and, in particular, Abelian
theorems, are the basis of the transcendental construc-
tion of the Jacobi variety of a Riemann surface. The
question of the inversion of an Abelian integral as a
function of its upper bound also leads to the concepts
of an Abelian function; an elliptic function; and theta-
functions (cf. Theta-function; Jacobi inversion problem).
Historically, the theory of Abelian integrals followed
from the consideration of a surface of genus g = 1. If
one writes the corresponding equation in the form
F(z, w) _ w 2 - fez) = 0,
15
ABELIAN INTEGRAL
where f (z) is a polynomial in z of the third or fourth
degree, then one obtains elliptic integrals (d. Elliptic
integral) as the respective Abelian integrals. They first
appeared at the end of the 17-th and the beginning of
the 18-th century as the result of the rectification of
curves of the second order in the studies of Jacob and
Johann Bernoulli and of G. Fagnano. L. Euler tackled
the addition theorem of elliptic integrals, which is a
special case of a theorem of N.H. Abel (1752). Abel
and e.G.J. Jacobi (1827) stated the problem of inver-
sion of elliptic integrals and obtained the solution. The
beginnings of the theory of elliptic functions were thus
established. However, some facts concerning this theory
had been established by e.F. Gauss early in the 18-th
century. Abel and Jacobi dealt with the much more dif-
ficult case of inversion of Abelian integrals in the case
g> 1. During the very first stages of development stress
was laid on hyper-elliptic integrals, where
F(z, w)=w2 - f(z) with f(z) a polynomial of the fifth
or sixth degree without multiple roots. Here g = 2 and
the difficulty of the inversion problem can already be
noticed. The principal advances in the theory of inver-
sion of Abelian integrals are due to B. Riemann (1851),
who introduced the concept of Riemann surfaces and
formulated and gave proofs of a large number of
important results.
Multi-dimensional generalizations of the theory of
Abelian integrals form the subject matter of algebraic
geometry and the theory of complex manifolds.
References
[I] CHEBOTAREV, N.G.: The theory of algebraic functions,
Moscow-Leningrad, 1948, Chapt. 8;9 (in Russian).
[2] SPRINGER, G.: Introduction to Riemann surfaces, Addison-
Wesley, 1957, Chapt.lO.
[3] NEVANLINNA, R.: Unijormisierung, Springer, 1953, Chapt 5.
[4] BLISS, G.A.: Algebraic functions, Amer. Math. Soc., 1933.
[5] STAHL, H.: Theorie der Abelschen Funktionen, Leipzig, 1896.
[6] SHAFAREVICH, I.R.: Basic algebraic geometry, Springer, 1977
(translated from the Russian).
E.D. Solomentsev
Editorial comments. An interesting and useful additional
reference is [A 1].
References
[A 1] LANG, S.: Introduction to algebraic and abelian functions,
Addison-Wesley, 1972.
AMS 1980 Subject Classification: 30F30, 26C15,
26E99
ABELIAN SCHEME - A smooth group scheme over
a base scheme S, the fibres of which are Abelian
varieties (cf. Abelian variety). The following is an
equivalent definition: An Abelian scheme over S, or an
Abelian S-scheme, is a proper smooth group S-scheme
all fibres of which are geometrically connected. Intui-
tively, an Abelian S-scheme may be understood as a
family of Abelian varieties parametrized by the scheme
16
S. A number of fundamental properties of Abelian
varieties carry over to Abelian schemes. For instance,
an Abelian S-scheme A is a commutative group S-
scheme [1], and if S is a normal scheme, A is projective
over S, [2].
Abelian schemes are used in the context of moduli
schemes of Abelian varieties with various auxiliary
structures, and also in the theory of reduction of
Abelian varieties (d. Neron model).
References
[1] MUMFORD, D.: Geometric invariant theory, Springer, 1965.
[2] RAYNAUD, M.: Faisceaux ampIes sur les schemas en groupes et
les espaces homogenes, Springer, 1970.
I. V Dolgachev
AMS 1980 Subject Classification: 14KXX
ABELIAN VARIETY - An algebraic group that is a
complete algebraic variety. The completeness condition
implies severe restrictions on an Abelian variety. Thus,
an Abelian variety can be imbedded as a closed sub-
variety in a projective space; each rational mapping of
a non-singular variety into an Abelian variety is regu-
lar; the group law on an Abelian variety is commuta-
tive.
The theory of Abelian varieties over the field of com-
plex numbers e is, in essence, equivalent to the theory
of Abelian functions founded by e.G.J. Jacobi, N.H.
Abel and B. Riemann. If en denotes n-dimension vec-
tor space, r c en is a lattice (cf. Discrete subgroup) of
rank 2n, then the quotient group X = en / r is a com-
plex torus. Meromorphic functions on X are the same
thing as meromorphic functions on en that are invari-
ant with respect to the period lattice r. If the field K
of meromorphic functions on X has transcendence
degree n, then X can be given the structure of an alge-
braic group. This structure is unique by virtue of the
compactness of X, and it is such that the field of
rational functions of this structure coincides with K.
The algebraic groups formed in this way are Abelian
varieties, and each Abelian variety over the field e
arises in this way. The matrix which defines a basis of
r can be reduced to the form (E IZ), where E is the
identity matrix and Z is a matrix of order n X n. The
complex torus X = en / r is an Abelian variety if and
only if Z is symmetric and has positive-definite ima-
ginary part. It should be pointed out that, as real Lie
groups, all varieties X are isomorphic, but this is not
true for their analytic or algebraic structures, which
vary strongly when deforming the lattice r. Inspection
of the period matrix Z shows that its variation has an
analytic character, which results in the construction of
the moduli variety of all Abelian varieties of given
dimension n. The dimension of the moduli variety is
n(n + 1) / 2 (cf. Moduli problem).

The theory of Abelian vanetIes over an arbitrary
field k is due to A. Wei! [1], [2]. It has numerous appli-
cations both in algebraic geometry itself and in other
fields of mathematics, particularly in number theory
and in the theory of automorphic functions. To each
complete algebraic variety, Abelian varieties (cf.
Albanese variety; Picard variety; Intennediate Jacobian)
can be functorially assigned. These constructions are
powerful tools in studying the geometric structures of
algebraic varieties. E.g., they were used to obtain one of
the solutions of the Liiroth problem. Another applica-
tion is the proof of the Riemann hypothesis for alge-
braic curves over a finite field - the problem for which
the abstract theory of Abelian varieties was originally
developed. It was also one of the sources of I-adic
cohomology. The simplest example of such a cohomol-
ogy is the Tate module of an Abelian variety. It is the
projective limit, as n-'HXl, of the groups X[ln] of points
of order In. The determination of the structure of such
groups was one of the principal achievements of the
theory of Wei!. In fact, if m is coprime with the charac-
teristic p of the field k and if k is algebraically closed,
then the group X[ m] is isomorphic to (Z / mZ)2dirnX. If
m =p, the situation is more complicated, which resulted
in the appearance of concepts such as finite group
schemes, formal groups and p-divisible groups (cf. Fin-
ite group scheme; Fonnal group; p-divisible group). The
study of the action of endomorphisms of Abelian
varieties, in particular of the Frobenius endomorphism
on its Tate module, makes it possible to give a proof of
the Riemann hypothesis (for algebraic curves over finite
fields, ct. Riemann hypotheses) and is also the principal
instrument in the theory of complex multiplication of
Abelian varieties. Another circle of problems connected
with the Tate module consists of a study of the action
of the Galois group of the closure of the ground field
on this module. There resulted the Tate conjectures
and the theory of Tate - Honda, which describes
Abelian varieties over finite fields in terms of the Tate
module [5].
The study of Abelian varieties over local fields,
including p-adic fields, is proceeding at a fast rate. An
analogue of the above-mentioned representation of
Abelian varieties as a quotient space en / r, usually
known as uniformization, over such fields, was con-
structed by D. Mumford and M. Raynaud. Unlike the
complex case, not all Abelian varieties, but only those
having a reduction to a multiplicative group modulo p,
are uniformizable [6]. The theory of Abelian varieties
over global (number and function) fields plays an
important role in Diophantine geometry. Its principal
result is the Mordell- Wei! theorem: The group of
rational points of an Abelian variety, defined over a
finite extension of the field of rational numbers, is fin-
itely generated.
ABNORMAL SUBGROUP
References
[1] WElL, A.: Courbes algebriques et varietes abeliennes. Varietes
abeliennes et courbes algebriques, Hermann, 1971.
[2] WElL, A.: Courbes algebriques et varietes abeliennes. Sur les
courbes algebriques et les varietes qui s'en deduisent, Hermann,
1948.
[3] WElL, A.: Introduction a /'etude des varietes kahteriennes, Her-
mann,1958.
[4] LANG, S.: Abelian varieties, Springer, 1983.
[5] MUMFORD, D.: Abelian varieties, Oxford Univ. Press, 1974.
[6] MANIN, Yu.I.: }I-Adie automorphic functions', J. Soviet Math.
5, no. 3 (1976), 279-333. (Itog; Nauk. i Tekhn. Sovrem Prob-
lemy 3 (1974), 5-93)
[7] SERRE, J.-P.: Groupes algebrique e{ corps des classes, Hermann,
1959.
[8] SIEGEL, C.L.: Automorphe Funktionen in mehrerer Variablen,
Math. Inst. Gottingen, 1955.
B.B. Venkov
A.N. Parshin
Editorial comments. For recent information on the Tate
conjectures see [A3]. For the theory of Tate-Honda see
also [A4]. Mumford's theory of uniformization is developed
in [A 1], [A2].
References
[A 1] MUMFORD, D.: 'An analytic construction of degenerating
curves over complete local rings', Compos. Math. 24 (1972),
129-174.
[A2] MUMFORD, D.: 'An analytic construction of degenerating
Abelian varieties over complete local rings', Compos. Math.
24 (1972), 239-272.
[A3] FALTINGS, G.: 'Endlichkeitssatze fi.ir abelsche Varietaten uber
Zahlkorpern', Invent. Math. 73 (1983), 349-366. (Errratum:
Invent. Math. 75 (1984), p. 381)
[A4] TATE, J.: 'Classes d'isoganie des variates aooliennes sur un
corps fini (d' apres T. Honda)', in sem. Bourbaki, Vol. 21,
1968-1969.
AMS 1980 Subject Classification: 14KXX
ABNORMAL SUBGROUP - A subgroup A of a
group G such that g E <A, A g > for any element g E G.
Here <A, A g > is the subgroup generated by A and its
conjugate subgroup A g = gAg -1. As an example of an
abnormal subgroup of a finite group G one can take
the normalizer (ct. Nonnalizer of a subset) NG(P) of
any Sylow p-subgroup peG, and even any maximal
subgroup NeG which is not normal in G. In the
theory of finite solvable groups (cf. Solvable group),
where many important classes of subgroups are abnor-
mal, use is made of the concept of a subabnormal sub-
group A of a group G, which is defined by a series of
subgroups
A =A0 C A Ie .. C An = G,
where Ai is abnormal in Ai + 1, i =0, ... ,n-1.
References
[1] HUPPERT, B.: Endliche Gruppen, I, Springer, 1967.
A.!, Kostrikin
Editorial comments. Nowadays, Ag is mostly defined as
A9=g-1Ag. Instead of 'solvable' also 'soluble' is frequently
used.
AMS 1980 Subject Classification: 20D25, 20D30,
20E07
17
ABSCISSA
ABsassA - One of the Cartesian coordinates of a
point.
AMS 1980 Subject Classification: 51 N 10
ABsOLUTE - l) The absolute of a regular topological
space X is the space aX which is mapped perfectly and
irreducibly onto X, and is such that any perfect irredu-
cible inverse image of the space X is homeomorphic to
aX. Each regular space X has a unique absolute. The
absolute of a space X is always extremally disconnected
and completely regular, and is perfectly and irreducibly
mapped onto X by means of a transformation
'lTx: aX~x. If two spaces X and Y are connected by a
single-valued or multi-valued perfect irreducible map-
ping f : X ~ Y, then their absolutes are homeomorphic,
and there exists a homeomorphism fa: aX~a Y such
that f='lTyfa'ITx 1.
If a homeomorphism fa: aX~a Y is given, the map-
ping f='lTyfa'ITX 1 will be, in the general case, multi-
valued, irreducible and perfect. Thus, absolutes and
their homeomorphisms 'control' the entire class of per-
fect irreducible mappings of regular spaces. The mean-
ing of this fundamental property is that absolutes of
regular topological spaces are projective objects in the
category of regular spaces and perfect irreducible map-
pings. If a regular space X is compact, respectively,
finally compact or complete in the sense of Cech, the
respective property is also displayed by the absolute of
this space. The absolute of a paracompact space is even
strongly paracompact, and is, moreover, perfectly zero-
dimensional. However, the absolute of a normal space
need not itself be normal. If X is a completely regular
space, then the Stone - Cech compactification of its
absolute is the absolute of any compactification of X.
Two spaces are called co-absolute if their absolutes are
homeomorphic.
Thus, the class of regular spaces is subdivided into
disjoint (pairwise not intersecting) classes of co-
absolute spaces. A space X is co-absolute with some
metric space if and only if it is a paracompact feath-
ered space containing a dense a-discrete system of open
sets. A compact space is co-absolute with some metriz-
able compactum if and only if it has a countable 'IT-
weight. If a compact space has a countable 'IT-weight
and has no isolated points, then and only then it will
be co-absolute with the perfect Cantor set. Conse-
quently, all metrizable compacta without isolated points
are co-absolute with the perfect Cantor set. The abso-
lute of a countable metrizable compactum is an exten-
sion of the Stone - Cech compactification of the natural
numbers. The absolute of an extremally disconnected
space is homeomorphic to it. Thus, the class of abso-
lutes (whatever this may be) of regular spaces coincides
with the class of extremally disconnected spaces. Since
lk
a non-discrete extremally disconnected space does not
contain any convergent sequence of pairwise distinct
points, the absolute of any non-discrete space is non-
metrizable (and does not even satisfy the first axiom of
countability).
Of the numerous ways in which the absolute aX of a
given (regular) space X can be constructed, the follow-
ing method is one of the simplest.
A family ~= {A} of non-empty canonical /Ca-sets, i.e.
of canonical closed sets A of the space X, is called a
thread if it is inclusion-directed, i.e. if for each two ele-
ments A, A' of the family ~ there exists an element A"
contained in A n A '. A thread ~ is called a maximal or
an end thread if it is not a subfamily of any thread dif-
ferent from it. It can be shown that threads exist;
moreover, it can be shown that for each non-empty /Ca-
set A the set DA of all threads containing the set A as
an element is non-empty. Each thread is contained in
some maximal thread. The intersection of all sets which
are elements of a maximal thread ~ is either empty or
contains a single point x(~); in the latter case the
thread ~ is said to be convergent (to the point x(~)). A
topology is introduced in the set aX of all ends (i.e.
maximal threads), by taking the collection of all sets
DA as a basis for the closed sets. The resulting topology
is Hausdorff and compact. The convergent ends in the
compactum aX form an everywhere-dense subspace.
The subspace of the space aX consisting of the conver-
gent ends is at the same time the absolute aX of X; it
turns out that aX is identical with the Stone-tech
compactification paX of ax' If X is not only regular,
but completely regular, the formula of commutativity of
the operators a and p is valid:
aX = paX = apX
v.I. Ponomarev
Editorial comments. The definition of the absolute of a
regular topological space X given above is slightly impre-
cise. A better (more precise) definition is: The absolute of a
regular topological space X is a pair (aX, 7Tx), where TTX is a
perfect irreducible mapping of aX onto X, such that for
every regular topological space Y and every perfect irredu-
cible mapping f of Yonto X there is a mapping g of aX onto
Y such that 7TX= fg.
In Western literature finally compact spaces are also
called Lindelof spaces and canonical closed set are also
called regular closed sets (or sometimes regularly closed
sets). The maximal or end threads are commonly called reg-
ular closed ultra filters.
The simple construction of ax given above can be
rephrased as follows: The family of regular closed sets
forms, in a natural way, a complete Boolean algebra. The
space aX is then simply the Stone space of this Boolean
algebra (the set of all ultrafilters (cf. Ultrafilter) on it, topolo-
gized using the sets OA as a base for the closed sets).
 ABSOLUTE CONTINUITY

References l0 such that the integral I if d/L I < for any set e
[A1J AlurnANGEL'SKIT, A.V. and PONOMAREV, V.I.: Fundamentals
of general topology. Problems and exercises, Reidel, 1984. with !L(e)<l>. In the general case the integral with
The O-absolute of a O-proximity space (X, l is the pair respect to a finitely-additive set function with scalar or
(X 8, '1Tx) consisting of a proximity space X8 and a pro- vectorial for /L is an absolutely continuous function.
jection '1Tx: X8~X which is a regular O-mapping. Here A.P. Terekhin
O-mapping is a term denoting any O-perfect, irreducible, VF. Emelyanov
O-proximity-continuous mapping. Any O-proximity 2) Absolute continuity of a measure as a concept in
space has a unique O-absolute. Any regular O-mapping the theory of measures. A measure p is absolutely con-
on a O-absolute is a proximity equivalence. The 0- tinuous with respect to a measure /L if p is an absolutely
absolute of a space (X, l is the maximal inverse image continuous set function with respect to /L. Thus, let p be
of the space (X, l under regular O-mappings. For each a finite measure, given together with /L on some fixed
regular O-mapping f: (X, l~(Y, l>') there exists a a-algebra G; p will then be absolutely continuous with
proximity equivalence F: X8~ Y 8' such that the respect to /L if it follows from !L(A) = 0, A E G, that
diagram F II{A ) = O. A generalized finite measure (d. Charge) p is
Xs ~ Ys' absolutely continuous with respect to a generalized
'1Tx~ ~'1Ty measure /L if II{A) = 0, provided that I /L I(A) = 0, where
X ~ Y I /L I is the total variation of /L. A.P. Terekhin
f
is commutative. 3) Absolute continuity of a function is a stronger
For maximal O-proximities on regular topological notion than continuity. A function f defined on a seg-
spaces the concept of a regular O-mapping is identical ment [a, b] is said to be absolutely continuous if for any
with that of a perfect irreducible mapping, while the >0 there exists a l0 such that for any finite system
concept of a O-absolute is identical with that of the of pairwise non-intersecting intervals (ak, bk)C(a, b),
absolute of a regular topological space. k = 1, ... , n, for which
n
References ~ (bk -ak) < 8,
[1) PONOMAREV, V.I.: 'On spaces co-abolute with metric spaces', k=i
Russian Math. Surveys 21, no. 4 (1966),87-114. (Uspekhi Mat.
Nauk 21, no. 4 (1966), 101-132)
the inequality
[2) GLEASON, A.M.: 'Projective topological spaces', Illinois J.
Math. 2, no. 4A (1958), 482-489. ~ Ij(bd-j(ad I<
k=i
[3) PONOMAREV, V.I.: 'Paracompacta, their projection spectra and
holds. Any absolutely continuous function on a seg-
continuous images', Mat. Sb. 60 (102), no. 1 (1963), 89-119 (in
Russian).
ment is continuous on this segment. The opposite
[4) FEDORCHUK, V.V.: 'Perfect irreducible mappings and general-
implication is not
ized proximities', Math. USSR-Sb. 5, no. 4 (1968), 498-508. true: e.g. the function
(Mat. Sb. 74 (118), no. 4 (1968),513-536) V V r:' d f(x)=x sin(l / x) if O<x~ 1 and f(O)=O is continu-
h k
. . re orc u
ous on the segment [0, 1], but is not absolutely continu-
2) The absolute in projective geometry is the curve
ous on it. If, in the definition of an absolutely continu-
(surface) of the second order constituting the set of
ous function, the requirement that the pairwise intersec-
infinitely-distant points in the Klein interpretation of a
tions of intervals (ab bk ) are empty be discarded, then
hyperbolic plane (space). The absolute can be used to
the function will satisfy an even stronger condition: A
introduce a metric on a projective plane (space) (d.
Lipschitz condition with some constant.
Projective detennination of a metric). For instance, the
If two functions f and g are absolutely continuous,
projective measure of a segment AB is defined as a
then their sum, difference and product are also abso-
quantity which is proportional to the natural logarithm
lutely continuous and, if g does not vanish, so is their
of the double ratio (ABCD) of four points, where C and
quotient f / g. The superposition of two absolutely con-
D are the points of intersection of the straight line AB
tinuous functions need not be absolutely continuous.
with the absolute.
A.B. Ivanov However, if the function f is absolutely continuous on a
AMS 1980 Subject Classification: 54C99, 51 N15 segment [a, b] and if A~f(x)~B, xE[a, b], while the
function F satisfies a Lipschitz condition on the seg-
ABSOLUTE CONTINUITY - 1) Absolute continuity of ment [A, B], then the composite function F[f(x)] is
an integral as a property of the (Lebesgue) integral. absolutely continuous on [a, b]. If a function f, which
Let a function f be /L-integrable on a set E. The integral is absolutely continuous on [a, b], is monotone increas-
of f over /L-measurable subsets e CE is an absolutely ing, while F is absolutely continuous on [f (a),f (b )],
continuous set function (see Subsection 3 below) with then the function F[f (x)] is also absolutely continuous
respect to the measure /L if for any >0 there exists a on [a, b].

19
ABSOLUTE CONTINUITY
An absolutely continuous function maps a set of
measure zero into a set of measure zero, and a measur-
able set into a measurable set. Any continuous function
of finite variation which maps each set of measure zero
into a set of measure zero is absolutely continuous.
Any absolutely continuous function can be represented
as the difference of two absolutely continuous non-
decreasing functions.
A function f that is absolutely continuous on the seg-
ment [a, b] has a finite variation on this segment and
has a finite derivative / (x) at almost every point. The
derivative / (x) is summable over this segment, and
x P(E)
I(x) = I(a) + jl(t)dt.
a
If the derivative of an absolutely continuous function is
almost everywhere equal to zero, then the function
itself is constant. On the other hand, for any function cp
that is summable on [a, b] the function f
cp(t)dt is
absolutely continuous on this segment. Accordingly, the
class of functions that are absolutely continuous on a
given segment coincides with the class of functions that
can be represented as an indefinite Lebesgue integral,
i.e. as a Lebesgue integral with a variable upper limit of
a certain summable function plus a constant.
If f is absolutely continuous on [a, b], then its total
variation is b Lebesgue (Haar) measure. A non-negative measure p.
VV(x) = j II (x) I dx.
a
The concept of absolute continuity can be general-
ized to include both functions of several variables and
set functions (see Subsection 4 below).
References
[I) KOI.MOGOROV, A.N. and FOMIN, S.V.: Elements of the theory
offunctions and functional analysis, Graylock, 19571961
(translated from the Russian).
[2) SMIRNOV, V.I.: A course of higher mathematics, 5, Addison-
Wesley, 1964 (translated from the Russian).
L.D. Kudryavtsev
4) Absolute continuity of a set function is a concept
usually applied to countably-additive functions defined
on a a-ring S of subsets of a set X. Thus, if p. and " are
two countably-additive functions defined on Shaving
values in the extended real number line [ - 00, 00], then
" is absolutely continuous with respect to p. (in symbols
this is written as "p.) if I p.1 (E)=O entails P(E)=O.
Here I p. I is the total variation of p.:
1p,1 =p,++p,-,
p,+ = sup{p.(F): E-:JFES},
p,- = -inf{p.(F): E~FES},
p.+ and p.- are measures, known as the positive and
negative variations of p.; according to the
Jordan - Hahn theorem, p. =p. + - p. -. It turns out that
20
the relations I) " p.; 2) ,,+ p., ,,- p.; 3)
I " IEO; I p. I are equivalent. If the measure " is finite,
"p. if and only if for any t:>0 there exists a 8>0
such that I p. I <8 entails I" I <t:o According to the
Radon-NikocJYm theorem, if p.," are (completely) 0'-
finite, (i.e. XES and there exists a sequence {En},
n=I,2, ... , such that
n91En = X, I p.(En) I, I P(En) I < 00 ]
and if "p., then there exists on X a finite measurable
function f such that
= ildp" EES.
Conversely, if p. is (completely) a-finite and the
I
integral dp. makes sense, then lei
dp. as a function of
the set E is absolutely continuous with respect to p.. If p.
and " are (completely) a-finite measures on (X, S);
there exist uniquely' defined (completely) a-finite meas-
ures "1 and "2 such that "="1 +"2, "1 p. and "2 is
singular with respect to p. (i.e. there exists a set A E S
such that 1"2 I(A)=O, I p.1 (X\A)=O) (Lebesgue's
theorem). A measure, defined on the Borel sets of a
finite-dimensional Euclidean space (or, more generally,
of a locally compact group), is called absolutely con-
tinuous if it is absolutely continuous with respect to the
on the Borel sets of the real line is absolutely continu-
ous if and only if the corresponding distribution func-
tion F(x)=p.{( - 00, x)} is absolutely continuous (as a
function of a real variable). The concept of absolute
continuity of a set function can also be defined for
finitely-additive functions and for functions with vector
values.
References
[I] HALMos, P.R.: Measure theory, v. Nostrand, 1950.
[2) NEVEU, J.: Bases mathematiques du calcul des probabilites, Mas-
~I~ V.V.S
. . azonov
Editorial comments.
References
[A 1] ROYDEN, H.L.: Real analysis, Macmillan, 1968.
[A2] ZAANEN, A.C.: Integration, North-Holland, 1967.
[A3] RUDIN, W.: Principles of mathematical analysis, McGraw-Hili,
1953.
[A4] RUDIN, W.: Real and complex analysis, McGraw-Hili, 1966.
[AS] TAYLOR, A.E.: General theory of functions and integration,
Blaisdell, 1965.
[A6] ALIPRANTZ, C.D. and BURLEINSHAW, 0.: Principles of real
analysis, North-Holland, 1981.
AMS 1980 Subject Classification: 26A 15, 28A 10
ABSOLUTE ERROR - See Error.
AMS 1980 Subject Classification: 65GXX
ABSOLUTE GEOMETRY - A geometry based on the
 ABSOLUTE SUMMABILITY

axioms of Euclidean geometry, except for the axiom on closed subset. A space satisfying the property in the article
parallel straight lines (postulate V). The assumptions above is then called an absolute extensor CAE) for normal
underlying absolute geometry are common to Euclidean spaces. One then proves that a space is an CAR) if and only
geometry and to Lobachevskii geometry. The term if it is an CAE). Absolute retracts and extensors can be
'absolute geometry' was introduced by J. Bolyai in defined for any class of spaces (i.e. not just for normal
spaces).
1832.
A.B. Ivanov A binormal space is a space X for which the product X x I
AMS 1980 Subject Classification: 51-01 is normal. It can be proved that a space X is binormal if and
only if X is normal and countably paracompact Ccf.
ABSOLUTE MOMENT of a random variable X - The Paracompactness criteria).
mathematical expectation of I X I r, r >0. It is usually References
[A1] Hu, S.T.: Theory of retracts, Wayne State Univ. Press,
denoted by fin so that Detroit, 1965.
f3, = EI X I' AMS 1980 Subject Classification: 54C55
The number r is called the order of the absolute
moment. If F(x) is the distribution function of X, then ABSOLUTE SUMMABILITY - A special type of
summability of series and sequences, which differs from
f3, = f
-xc
I x I' dF(x), (I) ordinary summability in that additional restrictions are
imposed. For matrix summation methods the require-
and, for example, if the distribution of X has density ment is that the series and sequences obtained as a
p(x), one has + CI.) result of the transformation corresponding to the given
f3, = f
-co
Ix I'p(x)dx. (2) summation method must be absolutely convergent (d.
Matrix summation method). Let a summation method A
In relation to the equations (I) and (2) one also be defined as a transformation of sequences {sn} into
speaks, respectively, of the absolute moments of the sequences {(In} by means of the matrix II ank II :
distribution function F(x) and the density p(x). The 00

existence of fir implies the existence of the absolute an = ~ankSb n=O, 1, ....
k=O
moment fir and also of the moments (d. Moment) of
order /. for 0</ ~r. Absolute moments often appear The sequence {sn} is then absolutely summable by the
in estimates of probability distributions and their method A (I A I-summable) to the limit s if it is A-
characteristic functions (d. Chebyshev inequality in pro- summable to this limit, i.e. if
bability theory; Lyapunov theorem). The function logfir lim an = s,
is a convex function of r, and the function fi~ / r is a n-->oo

non-decreasing function of r, r>O. and if the sequence {an} is of bounded variation:

Yu. V Prokhorov
AMS 1980 Subject Classification: 60A05, 60E05, (1)
n=l
60E10
If the Sn are the partial sums of a series
ABSOLUTE REfRACf FOR NORMAL SPACES - A
(2)
topological space X such that every mapping g: A ~ X
of any closed subset A of an arbitrary normal space Y then the series (2) is absolutely summable by the method
can be extended to the entire space Y. A direct product A to the sum s. Condition (1) is an additional require-
of absolute retracts for normal spaces is an absolute ment which makes absolute summability different from
retract for normal spaces, as is any retract (cf. Retract ordinary summability. Absolute summability is defined
of a topological space) of an absolute retract for normal in a similar manner for methods involving matrix
spaces. In particular, the following spaces are absolute transformations of series into sequences. If the summa-
retracts for normal spaces: the unit interval I; the n- tion method is defined by a transformation of the series
dimensional cube In; and the Hilbert cube IW. Any two (2) into a series 00

mappings of a binormal space into an absolute retract ~ an (3)

for normal spaces are homotopic; while a binormal n=O

absolute retract for normal spaces is contractible into a by means of a matrix II bnk II:
point. 00
M.l. Vottsekhovskit an = ~bnkUk'
k=O
Editorial comments. Usually an absolute retract CAR) for
normal spaces is defined to be a normal space which is a then the supplementary restriction is absolute conver-
retract of every normal space in which it is imbedded as a gence of the series (3). In the particular case in which

21
ABSOLUTE SUMMABILITY
the method A corresponds to the identity transforma-
tion of a sequence into a sequence (or of a series into a
series), absolute summability of a series is the same as
absolute convergence.
The supplementary requirements are suitably modi-
fied for non-matrix summation methods. Thus, in the
case of the Abel summation method such a requirement
is that the function
00
I(x) = ~UkXk
k=O
is of bounded variation on the semi-interval O:S;;;x < 1.
In the case of integral summation methods, absolute
summability is distinguished by the requirement of
absolute convergence of the correspondiong integrals.
Thus, for the Borel summation method the integral
00 00 U Xk
[e-xk~o-h-dx
must be absolutely convergent.
A summation method is said to preserve absolute con-
vergence of a series if it absolutely sums each absolutely
convergent series. If each such series is summable by
this method to a sum equal to the sum to which it con-
verges, then the method is called absolutely regular. For
instance, the Cesaro summation methods (C, k) are
absolutely regular for k ~O. The Abel method is abso-
lutely regular. The following are necessary and suffi-
cient conditions for the absolute regularity of a summa-
tion method defined by a transformation of a series
into a series by means of a matrix I bnk I :
00 00
~lbnkl';;;M, ~bnk=l,k=O,I, ...
n=O n=O
(the Knopp- Lorentz theorem). Summation methods by
means of other types of transformations are subject to
analogous restrictions.
A generalization of absolute summability is absolute
summability of degree p where p ~ 1. Here the condition
~ p-ll IJn-IJn-1 IP <
00
...:.,n 00.
n:=:: 1
is the supplementary restriction by which absolute sum-
mability of degree p is distinguished from ordinary
summability in the case of, say, a summation method
defined by a transformation of a sequence {sn} into a
sequence {on}.
The concept of absolute summability was introduced
by E. Borel for one of his methods, but was stated dif-
ferently from its modem formulation: absolute summa-
bility was distinguished by imposing the condition
eLf
oe - x
I
k~O
00 Uk'
/!xk I dx < 00
for each p =0, 1, .... Absolute summability was 101-
tially used in the study of the summability of power
series outside the disc of convergence. In view of the
22
problems involved in the multiplication of summable
series, absolute summability was defined and studied
by the method of Cesaro IC I, Ik I)-summability).
The general definition of absolute summability is more
recent, and has found extensive application in studies
on the summation of Fourier series.
References
[l] HARDy, G.H.: Divergent series, Clarendon, 1949.
[2] KOGBETLIANTZ, E.: Sommation des series et integrales diver-
gentes par les moyennes arithmetiques et typiques, Gauthier-
Villars, 1931.
[3] KNopp, K. and LORENTZ, G.G.: 'Beitriige zur absoluten Limi-
tierung', Arch. Math. (Basel) 2 (1949-1950), 10-16.
[4] KANGRO, G.F.: 'Summability theory of series and sequences',
ftogi Nauk. i Tekhn. Mat. Anal. 12 (1974),5-70 (in Russian).
1.1. Volkov
Editorial comments.
References
[A1] ZELLER, W. and BEEKMANN, W.: Theorie der Limitierungsver-
tahren, Springer, 1970.
AMS 1980 Subject Classification: 40F05
ABSOLUTE TOPOLOGICAL PROPERTY - An
obsolete term denoting the properties of a given set as
a topological space in contrast to properties related to
an imbedding of it in other spaces.
P.s. Aleksandrov
AMS 1980 Subject Classification: 54A05
ABSOLUTE VALUE, modulus, of a real number a -
The non-negative number, denoted by Ia I, which is
defined as follows: If a~O, lal=a; if a<O,
I a I = -a. The absolute value (modulus) of a complex
number z =x +ry, where x and yare real numbers, is
the number + Vx 2 +y2 . Absolute values obey the fol-
lowing relations
Ia I = I -a I, Ia 12 = Ia 2 = a 2 (a real),
1
lal-Ibl"';; la+bl",;; lal+lbl,
lal-Ibl"';; la-bl"';; lal+lbl.
labl = lal'lbl, andifb*o* = 1*1.
A generalization of the concept of the absolute value
to the case of arbitrary fields exists, cf. Norm on a field.
AMS 1980 Subject Classification: 26A03
ABSOLUTELY CONVERGENT IMPROPER
INTEGRAL - An improper integral such that the
integral of the absolute value of the integrand con-
verges. If an improper integral is absolutely convergent,
it is also convergent. To take a concrete example, let an
improper integral be given by
h ~
fl(x)dx = Iimfl(x)dx, -oo<a<b",;;+oo, (*)
(J l}-h {1
 ABSOLUTEL Y CONVERGENT SERIES

where the function f (x) is Riemann- (or Lebesgue-)
integrable on all intervals [a, 1]], a';;;;1]<b.
A necessary and sufficient condition for the absolute
convergence of the integral (*) (Cauchy's criterion for
the absolute convergence of an improper integral) is
that for any >0 there must exist an 1]0 a ';;;;1]( <b, such
that for any 1]' and 1]", 1]( ';;;;1]' <b, 1]( ';;;;1]" <b, the ine-
quality
The absolute convergence of the integral then implies
its convergence, but the converse proposition is not
true.
References
[I] IL'IN, V.A. and POZNYAK, E.G.: Foundations of mathematical
analysis, 1, Moscow, 1971 (in Russian).
[2] KUDRYAVTSEV, L.D.: Mathematical analysis, 1, Moscow, 1973
(in Russian).
[3] NIKOL'sKIi, S.M.: A course of mathematical analysis, 1, Mos-
cow, 1973 (in Russian).
L.D. Kudryavtsev

is true. Editorial comments.

If the improper integral is absolutely convergent, References
then it is equal to the Lebesgue integral of the (A1] ApOSTOL, T.M.: Mathematical analysis, Addison-Wesley,
1969.
integrand. There exist improper integrals which are
convergent but not absolutely convergent, for example AMS 1980 Subject Classification: 40A 10, 26A42
+00

j SlllX dx. ABSOLUTELY CONVERGENT SERIES - A series

o x

In order to find out whether or not a given integral is n~1

absolutely convergent it is expedient to use the tests for with (in general) complex terms for which the series
convergence of improper integrals of non-negative func- 00

tions; for example, the absolute convergence of ~ IUn I (2)

f
n=i

sin\lnx dx converges.
i x For a series (1) to converge absolutely, a necessary
is established with the aid of the comparison criterion and sufficient condition (Cauchy's criterion for the
of convergence. absolute convergence of a series) is that for any >0
For most of the available definitions of multiple there exist a number n( such that for all numbers
improper integrals there is a different relation between n >n( and for all integers p;;;'O the inequality
convergence and absolute convergence. Let a function
nfluk I< (
f (x) be defined on an open set G in the n-dimensional k=n
Euclidean space. If, for any sequence of cube-filled is true. If a series is absolutely convergent, it is also
domains Gk , k = 1, 2, ... , which monotonically convergent. The series
exhaust G (i.e. Gk C Gk C Gk + I, and
00 ~4,
U Gk = G), n=ln
k=i
where i = v=T, is absolutely convergent; the series
the limit of the Riemann integrals 00 i=JL
~
jf(x)dx, n=1 n
G,
is convergent, but not absolutely convergent. Let
as k~oo, exists and if this limit is independent of the
choice of the above sequence of domains, then it is usu- (3)
ally called the improper integral
be a series consisting of the same terms as the series (1)
jf(x)dx. but taken, in general, in a different order. Since the
G
series (1) is absolutely convergent, so is the series (3),
The integral thus defined is convergent if and only if it and the sum of the series (3) is equal to the sum of the
is absolutely convergent. There are also other defini- series (I). If the series
tions of multiple improper integrals. For example, for a 00 00

function f (x) which is defined on the whole space En ~un and ~vn
n= I n= I
and which is Riemann-integrable on any n-dimensional
sphere Q, of radius r < + 00 with centre at the origin, are both absolutely convergent, so is anyone of their
one can define the improper integral over En by the linear combinations
equation:
jf(x)dx = r~ If(X)dx.

23
ABSOLUTELY CONVERGENT SERIES
The series of all possible pairwise products Urn Vn of the
terms of these series, arranged in an arbitrary order, is
also absolutely convergent, and its sum is equal to the
product of the sums of the original series. These pro-
perties of absolutely convergent series are also
displayed by multiple series:
~ Un!> ... ,nk' (4)
(nJ, ... n,)
If a multiple series is absolutely convergent, it is con-
vergent, for example, both in the sense of spherical and
of rectangular partial sums, and its sums will be the
same in both cases. If the multiple series (4) is abso-
lutely convergent, the iterated series
(5)
nt=l nl=l
is absolutely convergent, i.e. all series obtained by suc-
cessive summation of terms of the series (4) by the
indices n 1, . . . ,nk are absolutely convergent; more-
over, the sums of the multiple series (4) and the iterated
series (5) are identical with the sum of any simple series
formed by all terms of the series (4).
If the terms of the series (1) are elements of some
Banach space with norm II . II, the series (1) is said to
be absolutely convergent if the series
n=i
is absolutely convergent. The properties of absolutely
convergent series of numbers discussed above can also
be generalized to include the case of absolutely
convergent series of elements of a Banach space. In
particular, (the partial sums of) an absolutely
convergent series in a Banach space converge(s) in that
space. The concept of an absolutely convergent series is
applied in a similar manner to multiple series in a
Banach space.
References
[I] IL'IN, V.A. and POZNYAK, E.G.: Foundations of mathematical
analySiS, I, Moscow, 1971 (in Russian).
[2] KUDRYAVTSEV, L.D.: Mathematical analySiS, I, Moscow, 1973
(in Russian).
[3] NIKOL'SKlI, S.M.: A course of mathematical analySiS, I, Mos-
cow, 1973 (in Russian).
L.D. Kudryavtsev
Editorial comments. A useful Western reference is [A 1].
References
[A1] ApOSTOL, T.M.: Mathematical analysis, Addison-Wesley,
1969.
AMS 1980 Subject Classification: 40A05
ABSOLUTELY-FLAT RING - A ring over which any
(right or left) module is a flat module. This class is
identical with the class of von Neumann regular rings
(d. Regular ring (in the sense of von Neumann.
AMS 1980 Subject Classification: 13C11, 16A30
24
ABSOLUTELY INTEGRABLE FUNCflON - A func-
tion such that its absolute value is integrable. If a
function is Riemann-integrable on the segment [a, b],
a <b, its absolute value is also Riemann-integrable on
this interval, and
[I I
I[f(X)dX I,,;; f(x) dx.
A similar assertion is also valid for a function of n vari-
ables which is Riemann-integrable on a cube-filled
domain in the n-dimensional Euclidean space. The con-
verse theorem for Riemann-integrable functions is not
true: The function f which is equal to I for rational
values of x and is equal to - I for irrational values of x
is not Riemann-integrable, but its absolute value is
integrable. The situation is different for Lebesgue-
integrable functions: A Lebesgue-measurable function f
is Lebesgue-integrable (Lebesgue-summable) on a
measurable set E in the n-dimensional space if and only
if its absolute value is Lebesgue-integrable on this set.
The following inequality is valid:
I!(x)dx I ,,;; I f(x) I dx.
In the case of improper one-dimensional Riemann or
Lebesgue integrals on a half-open interval [a, b),
a <b,,;;; + 00 (provided that the function f is, respec-
tively, Riemann-integrable or Lebesgue-integrable on
any segment [a,1)], a<1)<b), the existence of the
improper integral of the absolute value of the function,
b
j I f(x) I dx,
a
implies the existence of the integral
b
jf(x)dx,
a
but the converse statement is not true (d. Absolutely
convergent improper integral). In this connection it
should be noted that, if the improper integral
b ~
j I f(x) I dx = lim j I f(x) I dx
a 'r1-b a
exists, then f is Lebesgue-integrable on [a, b), and its
improper integral is equal to the Lebesgue integral.
In the case of functions of several (more than one)
variables, improper integrals are usually so defined that
the existence of the improper integral of the absolute
value of a function is equivalent to the existence of the
improper integral of the function itself.
Let the values of a function f be in some Banach
space with norm II II. The function f is then called
absolutely integrable on a measurable set E if the
integral
jllf(x)lldx
E
 ABSORPTION LAWS

exists; also, if f is integrable on E, the relationship An example of a Markov chain with absorbing state 0
is a branching process.
II [f(x)dX II.;;; [II f(x) II dx The introduction of additional absorbing states is a
convenient technique that enables one to examine the
is true. properties of trajectories of a Markov chain that are
References associated with hitting some set.
[1] IL'IN, V.A. and POZNYAK, E.G.: Foundations of mathematical Example. Consider the set S of states of a homo-
analysis, 1, Moscow, 1971 (in Russian).
[2] KUDRYAVTSEV, L.D.: Mathematical analysis, I, Moscow, 1973
geneous Markov chain ./) with discrete time and tran-
(in Russian). sition probabilities
[3] NIKOL'sKII, S.M.: A course of mathematical analysis, 1, Mos-
cow, 1973 (in Russian). Pi} = P{~t+1)=jl~t)=i},
[4] ScHWARTZ, L.: Cours d'analyse, I, Hennann, 1967. in which a subset H is distinguished and suppose one
L.D. KudryavIsev has to find the probabilities
Editorial comments. qih = P{7(H=h IO)=i}, iES, hEH,
References
[A1) ROYDEN, H.L.: Real analysis, Macmillan, 1968.
where T(H) = min{1 >0: T(t)EH} is the moment of first
hitting the set H. If one introduces the auxiliary Mar-
r
[A2) ZAANEN, A.C: Integration, North-Holland, 1967.
[A3) RUDIN, W.: Principles of mathematical analysis, McGraw-Hili, kov chain (I) differing from ./) only in that all states
1953.
[A4) RUDIN, W.: Real and complex analysis, McGraw-HIli, 1966.
r
h EH are absorbing in (I), then for h EH the proba-
[A5) TAYLOR, A.E.: General theory of functions and integration, bilities
Blaisdell, 1965.
[A6) ALIPRANTZ, C.D. and BURLEINSHAW, 0.: Principles of real = P{'T(H)';;;t, ~'T(H=h I~O)=i}
analysis, North-Holland, 1981.
are monotonically non-decreasing for It 00 and
AMS 1980 Subject Classification: 26A39, 26A42
qih = limp~h(t). iES, hEH. (*)
1-->00

ABSOLUTELY-UNBIASED SEQUENCE A By virtue of the basic definition of a Markov chain

sequence of random variables X\, ... , Xn for which
p~h(t+1) = llijP~h(t), t;;'O, iES\H, hEH,
the conditions JES

E(Xd = 0 and E(Xn+! I Xl> ... ,Xn) = 0 Phh(t) = 1, h EH; P~h(t) = 0, i, h EH, i=h.
are fulfilled, for n = 1, 2, .... The partial sums The passage to the limit for I~oo taking into acount
Sn =X\ + ... + Xn of an absolutely-unbiased sequence (*) gives a system of linear equations for qih:
form a martingale. These two types of sequences are qih = llUqih, iES\H, hEH,
interconnected as follows: The sequence {Yn } forms a jES

martingale if and only if it is of the form qhh = 1, hEH; qih = 0, i, hEH, i=h.
Yn =X\ + ... + Xn +c (n = 1, 2, ... , and c=E(Y\) is
a constant), where {Xn} is an absolutely-unbiased References
[1) W.FELLER: An introduction to probability theory and its applica-
sequence. Thus, all martingales are related to partial tions, 1, Wiley, 1968.
sums of certain absolutely-unbiased sequences. Simple A.M. Zubkov
examples of absolutely-unbiased sequences are AMS 1980 Subject Classification: 60J10
sequences of independent random variables with
mathematical expectation zero. Besides the term ABSORPTION LAWS - Identities of the form
'unbiased' the term 'fair' - with the related concept of xA(xVy) = x, xV(xAy) = x,
a 'fair play', is also employed.
A. V. Prokhorov where 1\ and V are two-place operations on some set
Editorial comments. In [A 1] the term 'absolutely fair' is L. If these operations satisfy also the laws of commuta-
used instead of absolutely-unbiased. tivity and associativity, then the relation x ~y defined
References by the equivalence
[A 1] FELLER, W.: An introduction to probability theory and its x.;;;y ~-uVy=y (*)
applications, 2, Wiley, 1966, p. 210.
(or equivalently, by the equivalence x.;;;y~xl\y=x) is
AMS 1980 Subject Classification: 60G42
an order relation for which x 1\Y is the infimum of the
elements x and y, while x Vy is the supremum. On the
ABSORBING STATE of a Markov chain ./) - A state other hand, if the ordered set (L,';;;;) contains an
i such that
infimum x 1\y and a supremum x Vy for any pair of
P{t)=i I~s)=i} = 1 for any t;;'s. elements x and y, then for the operations V and 1\ the

25
ABSORPTION LAWS
laws of absorption, commutativity and associativity, as
well as the equivalence (*) apply.
References
[11 RAsIOWA, E. and SIKORSKI, R.: The mathematics of
metamathematics, Polska Akad. Nauk, 1963. VN G . h'
. TIS In
Editorial comments. Instead of absorption laws one also
uses the term absorptive laws, cf. [A1], Chapt. 2, Sect. 4.
References
[A 1] CoHN, P.M.: Universal algebra, Reidel, 1981.
AMS 1980 Subject Classification: 04A05, 08AXX
ABSTRACf ALGEBRAIC GEOMETRY - The branch
of algebraic geometry dealing with the general proper-
ties of algebraic varieties (cf. Algebraic variety) over
arbitrary fields and with schemes (cf. Scheme), which
are their generalizations. The first studies in abstract
algebraic geometry appeared as early as the nineteenth
century, but the main development of the subject dates
back to the nineteen-fifties, with the creation of the
general theory of schemes by A. Grothendieck. Interest
in algebraic geometry over arbitrary fields arose in the
context of number-theoretical problems and, in particu-
lar, in the theory of equations with two unknowns.
Very important in the development of abstract alge-
braic geometry was the introduction of the concept of
the zeta-function of an algebraic curve by E. Artin in
1924 (cf. Zeta-function in algebraic geometry), and the
proof of the analogue of the Riemann hypothesis for
elliptic curves by H. Hasse in 1933. The theory of alge-
braic curves over an arbitrary field of constants, which
was developed at the same time, was of fundamental
importance in this proof.
The general development of the theory of rings and
fields in the first two decades of the twentieth century
prepared the ground for a systematic development of
higher-dimensional algebraic geometry over arbitrary
fields. In his series of articles (1933 - 1938), B.L. van
der Waerden based abstract algebraic geometry on the
theory of polynomial ideals. In particular, he developed
the intersection theory on a non-singular projective
algebraic variety. The results of the studies conducted
in this field are summarized in [4].
It was noted in 1940 by A. Weil that a proof of the
Riemann hypothesis for algebraic curves of arbitrary
genus involves the theory of higher-dimensional
varieties over arbitrary fields. He accordingly developed
the theory of abstract algebraic varieties (not neces-
sarily projective) over an arbitrary ground field, the
theory of divisors, the intersection theory for such
varieties, and the general theory of Abelian varieties,
which had been previously studied only from the ana-
lytic point of view. The appearance in 1946 of the book
by Weil [9] made the theory of valuations and field
26
theory (Weil's 'general points' ('generic points')
language) to be generally accepted as a base of abstract
algebraic geometry for a long time thereafter.
In the early nineteen-fifties powerful methods of
commutative algebra [6], [8] were brought to bear on
abstract algebraic geometry. The subject was further
modified by J.-P. Serre in his study [7] on coherent
algebraic sheaves (cf. Coherent algebraic sheaf). The
ideas and methods of homological algebra were here
introduced into algebraic geometry for the first time.
The development of abstract algebraic geometry
paralleled that of the concept of an algebraic variety.
Following the definition of an abstract algebraic variety
given by Weil, various generalizations of this concept
were proposed, and the concept of a scheme proved to
be the most useful. A systematic exposition of these
ideas and the development of the theory of schemes
was initiated in 1960 by A. Grothendieck in a series of
memoirs [5], in which the language of functors and of
the theory of categories was introduced into abstract
algebraic geometry, and many classical constructions in
the field were radically transformed.
The rapid development of abstract algebraic
geometry was due to the recognition of the fact that it
is possible, in the framework of the theory of schemes,
to apply to the 'abstract case' practically all of the
known concepts of the classical complex case and, in
particular, the homology theory of complex-analytic
manifolds. An important role in the development of
abstract algebraic geometry was played by the Weil
conjecture (1947), which postulated the existence of a
cohomology theory in which the Lefschetz formula for
the number of fixed points of a mapping would be
valid, and which established an intimate connection of
this hypothesis with purely arithmetical problems of
algebraic varieties (cf. Zeta-function in algebraic
geometry).
The concept of a topologized category (a Grothen-
dieck topology) found numerous applications, the
development of which laid the foundations of new
branches of abstract algebraic geometry: the theory of
representable functors (cf. Representable functor), for-
mal geometry, Weil cohomology; K-theory, and the
theory of group schemes (cf. Group scheme). The ideas
and methods developed in this way had their influence
in many branches of mathematics (commutative alge-
bra, category theory, the theory of analytic spaces and
topology).
The recent generalization of the concept of an alge-
braic variety to that of an algebraic space, which was
proposed in the late nineteen-sixties, made it possible
to extend the framework of abstract algebraic
geometry, and to connect it even more closely with
other branches of algebraic geometry.

References
[I] GROTHENDIECK, A.: 'The cohomology theory of abstract alge-
braic varieties', in Proc. Internat. Math. Congress Edinburgh,
1958, Cambridge Univ. Press, 1960, pp. 103-118.
[2] DOLGACHEV, LV.: 'Abstract algebraic geometry', J. Soviet
Math. 2, no. 3 (1974), 264-303. (Itogi Nauk. i Tekhn. Algebra
Topol. Geom. 10 (1972),47-112)
[3] SERRE, J.-P.: Matematika 7, no. 5 (1963), 3-93.
[4] GROBNER, W.: Moderne algebraische Geometrie; die
Idealtheoretische Grundlagen, Springer, Wien, 1949.
[5] GROTHENDIECK, A. and DIEUDONNE, J.: 'Elements de
geometrie algebrique', Pub!. Math. IRES 4.
[6] SAMUEL, P.: Methodes d'algebre abstraite en geometrie
algebrique, Springer, 1955.
[7] SERRE, J.-P.: 'Faiseaux algebriques coherentes', Ann. of Math.
(2) 61, no. 2 (1955), 197-278.
[8] ZARISKI, 0.: 'Algebraic geometry. The fundamental ideas of
abstract algebraic geometry', in Proc. Internat. Congress
Mathematicians Cambridge, 1950, Vol. 2, Amer. Math. Soc.,
1952, pp. 77-89.
[9] WElL, A.: Foundations of algebraic geometry, Arner. Math. Soc.,
1946.
I. V Dolgachev
Editorial comments. The systematic development of the
concept of a topologized category and its applications was
carried out by Grothendieck, P. Deligne and many others
(cf. [A1], [A2]).
References
[A 1] DELIGNE, P.: 'La conjecture de Weil I', Publ. Math, IHES 43
(1974), 273-307,
[A2] DELIGNE, P.: 'La conjecture de Weilil', Publ Math, IHES 52
(1980), 237-252,
AMS 1980 Subject Classification: 14-XX
ABSTRACf ANALYTIC FUNCflON, analytic map-
ping of Banach spaces - A function f (x) acting from
some domain D of a Banach space X into a Banach
space Y that is differentiable according to Frechet every-
where in D, i.e. is such that for any point a ED there
exists a bounded linear operator 8f(a, .) from X into Y
for which the following relation is true:
II h ,,-I. "/(a +h)-l(a)-f)/(a, h)" = 0,
lim
II h 11--+0
where 11'11 denotes the norm on X or on Y; 8f(a, h) is
called the Frechet differential of f at a.
Another approach to the notion of an abstract ana-
lytic function is based on differentiability according to
Gateaux. A function f (x) from D into Y is weakly ana-
lytic in D, or differentiable according to Gateaux in D, if
for each continuous linear functional y' on Y and each
element hEX the complex function y'(f(x+~h)) is a
holomorphic function of the complex variable ~ in the
disc I ~ I<p(x, h), where p(x, h)=sup{ I ~ I:
x + ~h ED}. Any abstract analytic function in a domain
D is continuous and weakly analytic in D. The converse
proposition is also true, and the continuity condition
can be replaced by local boundedness or by continuity
according to Baire.
The term 'abstract analytic function' is sometimes
employed in a narrower sense, when it means a func-
ABSTRACT ANALYTIC FUNCTION
tion f (z) of a complex variable z with values in a
Banach space or even in a locally convex linear topo-
logical space Y. In such a case any weakly analytic
function f (z) in a domain D of the complex plane C is
an abstract analytic function. One can also say that a
function f (z) is an abstract analytic function in a
domain DeC if and only if f (z) is continuous in D
and if for any simple closed rectifiable contour LcD
the integralkf (z) dz vanishes. For an abstract analytic
function f (z) of a complex variable z the Cauchy for-
mula (cf. Cauchy integral) is valid.
Let f (x) be a weakly analytic function in a domain
D of a Banach space X. Then f(x + ~h), as a function of
the complex variable ~, has derivatives of all orders in
the domain D = {~: x +~h ED}, hEX, t!Iese derivatives
being abstract analytic functions from D into Y. If the
set {x +~h: I ~ I.;;;; 1} belongs to D, then
00 1
I(x +h) = ~ -, f)n/(x, h),
n=O n.
where the series converges in norm, and
f)n/(x, h) = ~/(x +~h)1
dt f=O
= ~ ( I(x+~h~-n-l d~.
2m Id=l
A function y =P(x) from X into Y is called a
polynomial with respect to the variable x of degree at
most m if, for all x, h EX and for all complex t one has
m
P(x+~h) = ~P.(x,hW',
.=0
where the functions P .(x, h) are independent of ~. The
degree of P (x) is exactly m if Pm(x, h );~:O. A power
series is a series of the form ~:=oPn(x) where Pn(x)
are homogeneous polynomials of degree n so that
Pn(ax)=il Pn(x), x EX, for all complex a. An arbitrary
weakly convergent power series ~:=oPn(x) in a
domain D converges in norm towards some weakly
analytic function f(x) in D, and Pn(x) = 8nf (0, x) / n!,
OED. A function f(x) is an abstract analytic function
if and only if it can be developed in a power series in a
neighbourhood of all points a ED
n=O
where all Pn(h) are continuous in X.
Many fundamental results in the classical theory of
analytic functions - such as the maximum-modulus
principle, the uniqueness theorems, the Vitali theorem,
the Liouville theorem, etc. - are applicable to abstract
analytic functions if suitable changes are introduced.
The set of all analytic functions in a domain D forms a
linear space.
27
ABSTRACT ANALYTIC FUNCTION
The notion of an abstract analytic function can be
generalized to wider classes of spaces X and Y, such as
locally convex topological spaces, Banach spaces over
an arbitrary complete valuation field, etc.
References
[1] HILLE, E. and PHILLIPS, R.: Functional analysis and semi
groups, Amer. Math. Soc., 1957.
[2] EDWARDS, R.E.: Functional analySiS: theory and applications.,
Holt, Rinehardt, Winston, 1965.
[3] ScHWARTZ, L.: COUTS d'analyse, 2, Hermann, 1967.
A.A. Danilevich
E.D. Solomentsev
AMS 1980 Subject Classification: 46G20, 58812,
58C10,32A30
ABSTRACflON, MATHEMATICAL - Abstraction in
mathematics, or mental abstraction, is a significant
component of the mental activity aimed at the formula-
tion of basic mathematical concepts. The most typical
abstractions in mathematics are 'pure' abstractions,
idealizations and their various multi-layered superposi-
tions (see [5]).
The mental act of 'pure' abstraction consists of the
fact that in a certain situation being considered, atten-
tion is fixed only on certain (essential) features of the
objects being considered and on the interrelationships
between them, with the exclusion of other properties
and relationships which are considered as irrelevant.
The result of such an act of abstraction, firmed in an
appropriate language, begins to play the part of a gen-
eral concept. A typical example of mathematical
abstraction of this kind is abstraction by identification.
The mental act of idealization means that, in a cer-
tain situation being considered, one's imagination gen-
erates a certain concept which becomes the object con-
sidered by one's consciousness. The properties imparted
to the concept by one's imagination are not only those
actually posessed by the initial objects as a result of the
act of 'pure' abstraction, but also other, imagined pro-
perties, which reflect the original properties of the ini-
tial objects in a modified manner, or even properties
altogether absent in reality. One of the most tradi-
tional mathematical idealizations is the abstraction of
actual infinity, which leads to the idea of an actual
infinity. This abstraction is the base of the set-theoretic
development of mathematics. Another traditional ideal-
ization - the abstraction of potential realizability -
leads to the idea of a potential infinity. This abstraction,
in conjunction with the refusal to use the abstraction of
actual infinity, forms the base of constructive founda-
tions of mathematics.
The nature of any mathematical theory is largely
determined by the nature of the mathematical abstrac-
tion on which the formulation of the fundamental con-
cepts of this theory is based. The analysis of such
28
abstractions is one of the principal tasks of the founda-
tions of mathematics. A careful consideration of the
problems related to this circle of questions resulted in
the recognition of the fundamental importance of the
following factors: 1) the evaluation of abstract objects
resulting from the superposition of far-reaching ideali-
zations requires the development of special means for
their understanding; such a development is a difficult
task, which forms the subject of a special discipline:
semantics; and 2) the logical apparatus which may be
applied to any given mathematical theory essentially
depends on the nature of the fundamental concepts of
this theory, and hence on the nature of the
mathematical abstractions accepted in the formulation
of these concepts. See Intuitionism; Constructive
mathematics.
Major contributions to the analysis of abstractions
used in mathematics were made by L.E.J. Brouwer [I],
H. Weyl [2], D. Hilbert [3], A.A. Markov [4], and oth-
ers.
References
[1] BROUWER, L.E.J.: 'De onbetrouwbaarheid der logische prin-
cipes', Tijdschrift voor Wijsbegeerte 2 (1908), 152-158.
[2] WEYL, H.: Das Kontinuum, Chelsea, reprint, 1973.
[3] HILBERT, D.: Grundlagen der Geometrie, Springer, 1913,
Appendix VIII.
[4] MARKov, A.A.: On the logic of constructive mathematics, Mos-
cow, 1972 (in Russian).
[5] SHANIN, N.A.: 'Constructive real numbers and constructive
function spaces', Trudy Mat. Inst. Steklov. 67 (1962), 15-294 (in
Russian).
N.M. Nagornyf
AMS 1980 Subject Classification: 03A05
ABSTRACflON BY IDENTIFICATION - A manner
of forming general abstract concepts in which, when
considering concrete initial objects, only those differ-
ences are considered that, for any reason, are relevant
to the given situation, while other differences that are
regarded as irrelevant are ignored. Initial objects differ-
ing from one another only in irrelevant aspects are con-
sidered as identical. Linguistically, abstraction by iden-
tification is manifested by the fact that two similar ini-
tial objects are identified with each other, and are spo-
ken of as a single object, after having been assigned
some suitable term. For instance, by identifying similar
letters (words, alphabets) with each other, one arrives
at the concept of an abstract letter, word or alphabet
(see [l]); by identifying equivalent fundamental
sequences of rational numbers one arrives at the con-
cept of real numbers (cf. Real number); by identifying
isomorphic groups with each other, one arrives at the
concept of an abstract group, etc.
See also Abstraction, mathematical; Equivalence; Iso-
morphism.
References
[I] MARKOV, A.A.: Theory of algorithms, Israel Progr. Sci. Trans!.,

1961 (translated from the Russian). Also: Trudy Mat. Inst.
Steklov. 42 (1954).
[2] MARKov, A.A.: On the logic of constructive mathematics, Mos-
cow, 1972 (in Russian).
N.M. Nagornyi
AMS 1980 Subject Classification: 03A05
ABSTRACfION OF ACfUAL INFINITY - A
mathematical idealization, related to a certain form of
the idea of infinity in mathematics - the idea of the
so-called actual infinity.
As applied to constructive processes of potentially
infinite duration (e.g. the generation of successive posi-
tive integers starting from zero), abstraction of actual
infinity consists of ignoring the fact that such processes
do not terminate in principle and in considering the
results of such processes on the assumption that they
have terminated, viz. that the sets of objects have been
generated. The resulting sets (objects) are then mentally
regarded as actual 'finished' objects. The application of
the abstraction of actual infinity to the above example
makes it possible to consider the set of all non-negative
integers - the natural sequence - as a mathematical
object.
Logically, acceptance of the abstraction of actual
infinity leads to the acceptance of the law of the
excluded middle as a logical principle.
Abstraction of actual infinity is particularly impor-
tant in the construction of mathematics on the base of
the general theory of sets, established by G. Cantor.
Abstraction of actual infinity, being a far-reaching
idealization, especially so when repeatedly employed in
conjunction with other idealizations, generates objects
whose 'tangibility' is indirect, as a result of which the
problems involved in understanding propositions con-
cerning such objects meet with certain difficulties.
Unrestricted use of abstraction of actual infinity in
mathematics as a legitimate method of generating
mathematical objects met with objections by a number
of mathematicians (L. Kronecker, c.F. Gauss, D. Hil-
bert, H. Weyl, and others). Positive programs of
mathematical construction based on abstraction of
potential realizability, without recourse to abstraction
of actual infinity, were proposed by L.E.J. Brouwer (see
Intuitionism) and by A.A. Markov (d. Constructive
mathematics).
See also Abstraction, mathematical.
N.M. Nagornyi
AMS 1980 Subject Classification: 03A05
ABSTRACfION OF POTENTIAL REALIZABILITY -
A mathematical idealization, related to a certain form
of the concept of infinity in mathematics - the idea of
a potential infinity.
As applied to constructive processes which can, in
ACCELERATION OF CONVERGENCE
principle, be indefinitely extended (e.g. the successive
generation of positive integers starting from zero), the
abstraction of potential realizability consists in ignoring
any possible spatial, temporal or material obstacles to
the realization of each successive step of the process,
and to consider each step as potentially realizable. The
application of the abstraction of potential realizability
to the example given above is tantamount to assuming
that a unit can be added to any natural number, that it
is possible to form the sum of any two natural
numbers, etc., but it does not imply the idea of the
existence of the natural sequence as an actual 'infinite
object'.
The acceptance of the abstraction of potential real-
izability logically leads to the principle of mathematical
induction.
Abstraction of potential realizability plays a special
role in constructive mathematics, in which propositions
concerning the existence of constructive objects that
satisfy given conditions are regarded as propositions on
the potential realizability of such objects.
See also Abstraction, mathematical.
N.M. Nagornyi
AMS 1980 Subject Classification: 03A05
ACCELERATION OF CONVERGENCE (for an itera-
tive method) - The construction of a modification of
the iterative method, by means of that method, that
possesses a greater rate of convergence. The accelera-
tion methods (acceleration processes) used are quite
varied (see [1] - [6]) and depend both on the problem to
be solved and on the type of iterative method. In those
cases where the iterative method can be considered as a
particular case of a class of iterative methods contain-
ing free iteration parameters, acceleration of conver-
gence can be reduced to the problem of the optimal
choice of these parameters. The optimization problem
can take various forms, and may reduce, for example,
to a transition from a method of simple iteration,
u n+ 1 = un-T(Lun-f), (1)
for the solution of a system of linear algebraic equa-
tions
Lu = j, L = L' > 0, (2)
either to the Richardson method with Chebyshev
parameters or to the method of conjugate gradients.
The rate of convergence of similar classical iterative
methods depends on the condition number v(L) of the
matrix L and can be fairly slow for large v(L). In such
situations, and particularly for the solution of systems
of grid equations, modifications of these methods are
often used, these can be defined by the fact that they
are used not for (2), but for an equivalent system
B-ILu = B-Ij,
29
ACCELERATION OF CONVERGENCE

where B =B* >0 is a specially selected operator (see [2] References

[A1] AxELSSON, 0.: 'Conjugate gradient type methods for unsym-
- [4]).
metric and inconsistent systems of linear equations', Linear
The operator B- 1L is self-adjoint and positive on algebra and its Appl. 34 (1980), 1-66.
some Euclidean space and the rate of convergence of [A2] GoLUB, G.H. and LooN, C.F. VAN: Matrix computations,
the modifications obtained depends on P(B -1 L). Simi- North Oxford Academic, 1983.
[A3] GoLUB, G.H. and VARGA, R.S.: 'Chebyshev semi-iterative
lar modifications are also used for more general prob- methods, successive over-relaxation methods, and second-
lems, including non-linear problems (see Non-linear order Richardson iterative methods. Part I, II', Numerical
equation, numerical methods). In order to realise these Math. 3 (1961),147-156; 157-168.
modifications it is important to be able to solve the sys- [A4] MANTEUFFEL, T.A.: 'The Tchebychev iteration for nonsym-
metric linear systems', Numerical Math. 28 (1977),
tems Bv = g effectively, since, for example, modification 307-327.
of (1) leads to the relation [A5] MEIJEllINK, J.A. and VORST, H.A. VAN DER: 'An iterative solu-
tion method for linear systems of which the coefficient matrix
(3) is a symmetric M-matrix', Math. Compo 31 (1977), 148-162.
(see Minimization of the labour of calculation). [A6] STEWART, G.W.: 'Conjugate gradients methods for solving
systems of linear equations', Numerical Math. 21 (1973),
One of the traditional general methods for the 284-297.
acceleration of convergence for methods (1) is the 82 _ [A7] VARGA, R.S.: Matrix iterative analysis, Prentice Hall, 1962.
process. It is used along with a whole series of other AMS 1980 Subject Classification: 65F10
acceleration methods (see [1]) in iterative methods for
the partial eigen value problem. ACCUMULATION OF ERRORS in the numerical
In solving non-linear problems, acceleration of con- solution of algebraiC equations - The overall effect of
vergence is often achieved by choosing a special initial rounding-off at the various stages of a computation
approximation on the basis of methods of continuation procedure on the accuracy of the computed solution to
by a parameter. For these same problems acceleration a system of algebraic equations. The most commonly
of convergence is also sometimes realized by the use of employed technique for a priori estimation of the total
iterative methods of a higher order (the effect of rounding-off errors in numerical methods of
Newton - Kantorovich method and others). linear algebra is the scheme of inverse (or backward)
Various methods for the acceleration of convergence analysis. Applied to the solution of linear algebraic
are also used in probabilistic iterative methods of equations
Ax = b (1 )
Monte-Carlo type (see [2]).
References the scheme of inverse analysis is as follows. On the
[I] FADDEEV, D.K. and FADDEEVA, V.N.: Computational methods assumption that some direct method M has been used,
of linear algebra, Freeman, 1963 (translated from the Russian). the computed solution XM does not satisfy (1), but it
[2] BAKHVALOV, N.S.: Numerical methods: analySis, algebra, ordi-
nary differential equations, Mir, 1977 (translated from the Rus- can be expressed as the exact solution of a perturbed
sian). system
(2)
[3] MARCHUK, GJ.: Methods of computational mathematics, Mos-
cow, 1980 (in Russian). The quality of the direct method is estimated in terms
[4] SAMARSKIl, A.A. and NIKOLAEV, E.S.: Methods of solving differ-
ence equatiOns, Moscow, 1978 (in Russian).
of the best a priori estimate that can be found for the
[5] GLOWINSKI, R.: Numerical methods for nonlinear variational norms of the matrix FM and the vector k M . These 'best'
problems, Springer, 1984. FM and kM are known as the equivalent perturbation
[6] D'YAKONOV, E.G.: 'On iterative methods with saddle opera- matrix and vector, respectively, for the method M.
tors', Dokl. Akad. Nauk. SSSR, New Series 292 (1987), 1037-
If estimates for FM and kM are available, the error of
1041 (in Russian). E.G. D'yakonov
the approximate solution XM can be estimated theoreti-
Editorial comments. In the Western literature, cally by the inequality
Richardson's method with Chebyshev parameters is usually
called the Chebyshev semi-iterative method. Chebyshev II X-XM II,,;::: cond(A) [II FM II + II kM II] 3
iteration is extensively analyzed in [A3] and [A7] , Chapt. 5. II x II "" I -II A -I II II FM II II A II IfbT I )
A further analysis for the case when L is non-symmetric can Here cond(A) = II A II II A -1 II is the condition number
be found in [M]. of the matrix A, and the matrix norm in (3) is assumed
Unlike Chebyshev iteration, the method of conjugate gra-
to be subordinate to the vector norm II II.
dients does not need accurate estimates of the largest and
In reality, an estimate for II A -1 II is rarely known
smallest eigen values of the iteration matrix 1- TL in (1).
Good references are [A2] , pp. 362-379, and [A6]. Further
in advance, and the principal meaning of (2) is the pos-
developments are discussed in [AS].
sibility that it offers of comparing the merits of dif-
The pre-multiplication by a matrix B- 1 of the system (2) ferent methods. In the sequel some typical estimates for
is often called pre-conditioning. An up-to-date survey of this the matrix FM are presented.
technique is given in [A 1]. For methods with orthogonal transformations and

30

floating-point arithmetic (A and b in the system (l) are
assumed to be real)
(4)
In this estimate, is the relative precision of in the
computer, II A liE = (~at)1 /2 is the Euclidean matrix
norm, and f(n) is a function of type Cnk, where n is
the order of the system. The exact values of the con-
stants C and the exponents k depend on such details of
the computation procedure as the rounding-off method,
the use made of accumulation of inner products, etc.
Most frequently, k = I or 3/2.
In Gauss-type methods, the right-hand side of the
estimate (4) involves yet another factor g(A), reflecting
the possibility that the elements of A may increase at
intermediate steps of the method in comparison with
their initial level (no such increase occurs in orthogonal
methods). In order to reduce g(A) one resorts to vari-
0us ways of pivoting, thus putting bounds on the
increase of the matrix elements.
In the square-root method (or Cholesky method),
which is commonly used when the matrix A is positive
definite, one has the sharper estimate
IIFMIIE';;; ellA IIE('
There exist direct methods (the methods of Gordan,
bordering and of conjugate gradients) for which a
direct application of a scheme of inverse analysis does
not yield effective estimates. In these cases different
arguments are utilized to investigate the accumulation
of errors (see [6] - [9]).
References
[I] GIVENS, W.: 'Numerical computation of the characteristic
values of a real symmetric matrix', Us. Atomic Energy Com
mission Reports. Ser. ORNL, no. 1574 (1954).
[2] WILKINSON, I.H.: Rounding errors in algebraic processes, Pren-
tice Hall, 1962.
[3] WILKINSON, I.H.: The algebraic eigenvalue problem, Oxford
Dniv. Press, 1969.
[4] VOEVODIN, V.V.: Rounding-off errors and stability in direct
methods of linear algebra, Moscow, 1969 (in Russian).
[5] VOEVODIN, V.V.: Computational fundamentals of linear algebra,
Moscow, 1977 (in Russian).
[6] PETERS, G. and WILKINSON, I.H.: 'On the stability of
Gauss-Gordan elimination with pivoting', Commun. ACM 18,
no. I (1975), 20-24.
[7] BROYDEN, CG.: 'Error propagation in numerical processes', J.
Inst. Math. and Appl. 14, no. 2 (1974), 131-140.
[8] REID, I.K.: Large sparse sets of linear equations, London New
[9]
York, 1971, pp. 231-254.
IKRAMov, KH.D.: 'The condition for intermediate matrices for
Gauss, Gordan and optimal elimination methods', Zh. Vychisl.
Mat. i Mat. Fiz. 18, no. 3 (1978), 531-545 (in Russian).
Kh.D.lkramov
Editorial comments. Reference [3] above contains a
detailed account of rounding-off error analysis (also called
roundoff error analysis) of algebraic processes. In addition
one can mention [A1], pp. 87-97, and [A5], pp. 69-82. A
historical analysis of roundoff analysis is given in [A7]. More
ACCUMULATION OF ERRORS
recent developments (interval analysis and automatic error
analysis) can be found in [A2], [A3] and [A8]. A complete
introduction to interval analYSis is given in [A4].
References
[A1] FORSYTHE, G.E. and MOLER, CB.: Computer solution of
linear algebraic systems, Prentice Hall, 1967.
[A2] LARSEN, I. and SAMEH, A.: 'Efficient calculation of the effects
of roundoff errors', ACM Trans. Math. Software 4 (1978),
228-236.
[A3] MILLER, W. and SPOONER, D.: 'Software for roundoff
analysis II', ACM Trans. Math. Software 4 (1978), 369-390.
[A4] STEWART, G.W.: 'Conjugate gradients methods for solving
systems of linear equations', Numer. Math. 21 (1973), 284-
297.
[AS] WILKINSON, I.H.: 'Error analysis of direct methods of matrix
inversion', J. Assoc. Comput. Machinery 8 (1961),281-330.
[A6] WILKINSON, I.H.: 'Modern error analysis', SIAM Review 13,
548-588.
[A7] YOHE, I.M.: 'Software for interval arithmetic', ACM Trans.
Math. Software 5 (1979), 50-63.
Accumulation of rounding-off errors (or of errors)
occurs when one is solving problems in which the solu-
tion is the result of a large number of consecutively
performed arithmetic operations.
A significant proportion of such problems arises in
the solution of linear or non-linear algebraic problems
(see above). In turn, one is most commonly concerned
with algebraic problems thai arise from the approxima-
tion of differential equations. These problems display
certain specific features.
Errors accumulate in accordance with the same - or
even simpler - laws as those governing the accumula-
tion of computational errors; they may be investigated
when analyzing a method for the solution of a problem.
There are two different approaches to investigating
the accumulation of computational errors. In the first
case one assumes that the computational errors at each
step are introduced in the most unfavourable way and
so obtains a majorizing estimate for the error. In the
second case, one assumes that the errors are random
and obey a certain distribution law.
The nature of the accumulation of errors depends on
the problem being solved, the method of solving and a
variety of other factors that appear at first sight to be
rather inessential: such as the type of computer arith-
metic (fixed-point or floating-point), the order in which
the arithmetic operations are performed, etc. For exam-
ple, in computing the sum of N numbers
AN = aj + ... +aN
the order of the operations is significant. Suppose that
the computation is being done in a computer with
floating-point arithmetic with t binary digits, all
numbers lying in the interval 1/2< I an I.;;; 1. In a
direct computation of AN via the recurrence formula
the majorizing error estimate is of the order 2- 1N. But
31
ACCUMULATION OF ERRORS
one can proceed otherwise (see [1]). First compute sums
of pairs, Al =a2k-1 +a2k (if N=21 + 1 is odd, one puts
AJ+I =a2/+1)' Then compute further sums of pairs,
A~=A!k-1 +A!k> etc. If 2m - I <N.;;;;,2m , then m steps
of pairwise addition using the formulas
AZ = A~-~l +A~-l, A2 _ ak>
yield AT =AN; the majorizing error estimate is of the
order 2-I logz N.
In typical problems, the numbers am are computed
according to formulas, in particular recurrence formu-
las, or appear consecutively in the operating memory of
the computer; in such cases use of the scheme just
described increases the load on the computer memory.
However, the sequence of computations can be so
organized that the load on the operating memory is
never more than ""' logz N cells.
In the numerical solution of differential equations
one may encounter the following cases. As the grid
spacing h tends to zero, the error increases as (a(h)i-',
where q >0, while limh-->o I a(h) I > 1. Such methods of
solution fall into the category of unstable methods. They
are of little practical interest.
In stable methods the error will characteristically
increase at the rate A (h)h -q, where
li~_o I A (h) I < 00. The error in such methods is usu-
ally estimated in the following way. One constructs an
equation for the perturbation, whether due to the intro-
duction of rounding-off errors or other errors, and then
investigates the solution of this equation (see [2], [3]).
In more complicated cases, one uses the method of
equivalent perturbations (see [1], [4]), developed in con-
nection with the investigation of the accumulation of
computational errors in the solution of differential
equations (see [3], [5], [6]). Computations according to
some standard scheme with rounding-off are treated as
computations without rounding-off but for equations
with perturbed coefficients. Comparing the solution of
the original difference equation with the solution of the
equation with perturbed coefficients, one obtains an
estimate for the error.
Considerable care is taken to select a method in
which q and A(h) are as small as possible. When a
fixed method of solution is being used, the standard
formulas can usually be transformed in such a way that
q';;; I (see [3], [5]). This is particularly important in the
case of ordinary differential equations, where the
number of steps in the various cases turns out to be
very high.
The magnitude A (h) can increase strongly as the
integration interval is enlarged. The goal is therefore to
use methods with A (h) as small as possible. When deal-
ing with Cauchy problems, the rounding-off error at
each specific step can be regarded, relative to the previ-
32
ous step, as an error in initial data. A lower bound for
A(h) is therefore dependent on the behaviour of the
difference between slightly different solutions of the
differential equation, as defined by the small-
disturbance equation.
In the numerical solution of an ordinary differential
equation y' = f(x, y), the small-disturbance equation
has the form ,
11 = hex, y)1! + S,
and so, when solving the problem on an interval
(Xo, X), one cannot reckon on the constant A(h) in a
majorizing estimate for the computational error being
essentially better than
x
X [h(t.Y(I))dl
fe dx.
Xo
When solving this problem, therefore, the most com-
monly used methods are one-step methods of the
Runge- Kutta or Adams type (see [3], [7]), where the
accumulation of errors is basically determined by the
solution of the small-disturbance equation.
In many methods, the principal term of the error
increases according to a similar law, while the computa-
tional error builds up significantly more rapidly (see
[3]). The domain of practical applicability for such
methods is substantially narrower.
The accumulation of computational errors essentially
depends on the method employed to solve the grid
problem. For example, when solving grid boundary
value problems corresponding to ordinary differential
equations by the shooting or double-sweep method, the
accumulation of errors behaves like A (h)h -q, where q
is the same. The magnitudes A (h) in these methods
may differ to such an extent that in a certain situation
one of the methods becomes useless. In the shooting
method for a grid boundary value problem, the accu-
mulation of errors behaves like c 1/ h, C >0, while for
the double-sweep method this is Ah -q. If one adopts a
probabilistic approach to the investigation of the accu-
mulation of errors, there are cases in which one
assumes a priori a certain distribution law for the
errors (see [2]); in others, one introduces a measure on
the space of problems under consideration and, on the
basis of this measure, derives a distribution law for the
rounding-off errors (see [8], [9]).
When measuring the accuracy of the solution of a
problem, the majorizing and probabilistic approaches
to the estimation of the accumulation of computational
errors usually yield qualitatively identical results: The
errors either accumulate in both cases within admissible
bounds, or they exceed them in either case.
References
[1) VOEVODIN, V.V.: Computational foundations of linear algebra,
Moscow. 1977 (in Russian).

[2] SHURA-BURA, M.R.: Prikl. Mat. i Mekh. 16, no. 5 (1952), 575-
588.
[3] BAKHVALOV, N.S.: Numerical methods: analysis, algebra, ordi-
nary differential equations, Mir, 1977 (translated from the Rus-
sian).
[4] WILKINSON,l.H.: The algebraic eigenvalue problem, Oxford
Univ. Press, 1969.
[5] BAKHVALOV, N.S.: Computational methods and programming,
Vol. I, Moscow, 1962, pp. 69-79 (in Russian).
[6] GoDUNOV, S.K. and RYABEN'KIi, V.S.: The theory of difference
schemes, North-Holland, 1964 (translated from the Russian).
[7] BAKHVALOV, N.S.: 'On an estimate of the error at numerical
integration of differential equations by Adams' extrapolation
method', Dokl. Akad. Nauk SSSR 104, no. 5 (1955), 683-686
(in Russian).
[8] BAKHVALOV, N.S.: 'Optimal convergence bounds for quadratic
processes and integral methods of Monte Carlo type for classes
of functions', Zh. Vychisl. Mat. i Mat. Fiz. 4, no. 3 (1964),
399-404 (in Russian).
[9] LAPSHIN, E.A.: 'A statistical investigation of roundoff errors in
the numerical solution of an ordinary differential equation',
Zh. Vychisl. Mat. i Mat. Fiz. 11, no. 6 (1971), 1425-1436 (in
Russian).
N. S. Bakhvalov
Editorial comments. Statistical treatment of rounding-off
(or roundoff) errors can be found in the references below. It
should be remarked that statistical analyses give more real-
istic bounds, but are much more difficult to compute.
References
[A 1] P.: Discrete variable methods in ordinaT'/ differential
HENRIe!,
equations, Wiley, 1962.
[A2] HENRIe!, P.: Error propagation for difference methods, Wiley,
1963.
[A3] HULL, T.E. and SWENSON, J.R.: 'Tests of probabilistic models
for propagation of roundoff errors', Commun. ACM 9 (1966),
108-113.
AMS 1980 Subject Classification: 65GXX, 65F35,
65LXX
ACCUMULATION POINT of a set A - A point x in a
topological space X such that in any neighbourhood of
x there is a point of A distinct from x. A set can have
many accumulation points; on the other hand, it can
have none. For example, any real number is an accu-
mulation point of the set of all rational numbers in the
ordinary topology. In a discrete space, no set has an
accumulation point. The set of all accumulation points
of a set A in a space X is called the derived set (of A).
In a T I-space, every neighbourhood of an accumula-
tion point of a set contains infinitely many points of
the set.
The concept just defined should be distinguished
from the concepts of a proximate point and a complete
accumulation point. In particular, any point of a set is a
proximate point of the set, while it need not be an
accumulation point (a counterexample: any point in a
discrete space).
A. V Arkhangel'skif
AMS 1980 Subject Classification: 54AXX
ACTION - A functional expressed by the definite
ACTION OF A GROUP ON A MANifOLD
integral of a function, the stationary values of which
determine a real motion of a mechanical system acted
upon by given active forces, in the class of kinemati-
cally possible motions between some two final positions
Po and P I in space and satisfying certain conditions.
One distinguishes between Hamiltonian, Lagrangian
and Jacobian actions, which appear in the correspond-
ing principles of stationary action.
The Hamiltonian action
1\
j(T+U)dt
10
is defined in the class of kinematically possible motions
of a holonomic system for which the initial and the
final positions of the system, as well as the time of
motion between them, are the same as the respective
ones for real motion.
The Lagrangian action
t
j2Tdt
10
and the Jacobian action
are defined in the class of kinematically possible
motions of a holonomic conservative system for which
the initial and the final positions of the system, and the
constant energy h, are the same as the respective mag-
nitudes for a real motion. Here T is the kinetic energy
of the system, and for a conservative system
1
T = -2 ~ a;/J;qj'
;.j=l
where qi are the generalized Lagrange coordinates, and
U(q) is the force function of active forces.
For more details see Variational principles of classical
mechanics; Hamilton - Ostrogradski principle; Lagrange
principle; Jacobi principle.
V V Rumyantsev
AMS 1980 Subject Classification: 70-01, 70BXX
ACTION OF A GROUP ON A MANIFOLD - The
best-studied case of the general concept of the action of
a group on a space. A topological group G acts on a
space X if to each g E G there corresponds a homeomor-
phism CPg of X (onto itself) satisfying the following con-
ditions: I) CPg 'CPh = CPgh; 2) for the unit element e E G the
mapping CPe is the identity homeomorphism; and 3) the
mapping cp: G XX~X, cp(g, x)=CPg(x) is continuous. If
X and G have supplementary structures, the actions of
G which are compatible with such structures are of spe-
cial interest; thus, if X is a differentiable manifold and
33
ACTION OF A GROUP ON A MANIFOLD
G is a Lie group, the mapping cp is usually assumed to
be differentiable.
The set {cpg(xo) }gEG is called the orbit (trajectory) of
the point Xo EX with respect to the group G; the orbit
space is denoted by X / G, and is also called the quo-
tient space of the space X with respect to the group G. An
important example is the case when X is a Lie group
and G is a subgroup; then X / G is the corresponding
homogeneous space. Classical examples include the
spheres S n -I = 0(n) / O(n - I), the Grassmann mani-
folds O(n) / (O(m)XO(n -m, and the Stiefel mani-
folds O(n) / O(m) (cf. Grassmann manifold; Stiefel
manifold). Here, the orbit space is a manifold. This is
usually not the case if the action of the group is not
free, e.g. if the set X G of fixed points is non-empty. A
free action of a group is an action for which g = e fol-
lows if gx =x for any x EX. On the contrary, X G is a
manifold if X is a differentiable manifold and the
action of G is differentiable; this statement is valid for
cohomology manifolds over Zp for G = Zp as well
(Smith's theorem).
If G is a non-compact group, the space X / G is usu-
ally inseparable, and this is why a study of individual
trajectories and their mutual locations is of interest.
The group G = R of real numbers acting on a differenti-
able manifold X in a differentiable manner is a classical
example. The study of such dynamical systems, which
in terms of local coordinates is equivalent to the study
of systems of ordinary differential equations, usually
involves analytical methods.
If G is a compact group, it is known that if X is a
manifold and if each gEG, g=l=e, acts non-trivially on
X (i.e. not according to the law (g, X)-H), then G is a
Lie group [8]. Accordingly, the main interest in the
action of a compact group is the action of a Lie group.
Let G be a compact Lie group and let X be a com-
pact cohomology manifold. The following results are
typical. A finite number of orbit types exists in X, and
the neighbourhoods of an orbit look like a direct pro-
duct (the slice theorem); the relations between the coho-
mology structures of the spaces X, X / G and X G are
of interest.
If G is a compact Lie group, X a differentiable mani-
fold and if the action
</>:GXX ~ X
is differentiable, then one naturaly obtains the follow-
ing equivalence relation: (X, CP)--(X', CP') if and only if
it is possible to find an (X", CP") such that the boundary
ax" has the form ax" =X U X' and such that
cp" I x = cp, cp" I x = cp'. If the group G acts freely, the
equivalence classes can be found from the one-to-one
correspondence with the bordisms Q.(B G ) of the classi-
fying space BG (cf. Bordism).
34
Recent results (mid-nineteen-seventies) mostly con-
cern: 1) the determination of types of orbits with vari-
ous suplementary assumptions concerning the group G
and the manifold X ([6]); 2) the classification of group
actions; and 3) finding connections between global
invariants of the manifold X and local properties of the
group actions of G in a neighbourhood of fixed points
of XG. In solving these problems an important part is
played by: methods of modem differential topology
(e.g. surgery methods); KG-theory [1], which is the
analogue of K-theory for G-vectorbundles; bordism and
cobordism theories [3]; and analytical methods of
studying the action of the group G based on the study
of pseudo-differential operators in G-bundles [2], [7].
References
[II ATIYAH, M.F.: K-theory: lectures, Benjamin, 1967.
[2) ATIYAH, M.F. and SINGER, I.M.: 'The index of elliptic opera-
tors', Ann. of Math. (2) 87 (1968), 484-530.
[31 BUKHSHTABER, V.M., MISHCHENKO, A.S. and NOVIIWV, S.P.:
'Formal groups and their role in the apparatus of algebraic
topology', Russian Math. Surveys 26 (1971), 63-90. (Uspekhi
Mat. Nauk 26, no. 2 (1971), 131-154)
[4) CONNER, P.E. and FLOYD, E.E.: Differentiable periodic maps,
Springer, 1964.
[5) BREDON, G.: Introduction to compact transformation groups,
Acad. Press, 1972.
[6) HSIANG, Wu YI: Cohomology theory of topological transforma-
tion groups, Springer, 1975.
[7) ZAGIER, D.B.: Equivariant Pontryagin classes and applications to
orbit spaces, Springer, 1972.
[8) Proc. conf. transformation groups, Springer, 1968.
[9) Proc. second conf. compact transformation groups, Springer,
1972.
A. V. Zarelua
Editorial comments.
References
[A1] PETRIE, T. and RANDALL, J.D.: Transformation groups on
manifolds, M. Dekker, 1984.
AMS 1980 Subject Classification: 57SXX
ACYCLIC CONTINUUM - A continuum with trivial
homology groups (d. Homology group).
P.s. Aleksandrov
AMS 1980 Subject Classification: 54F15
ACYCLIC ELEMENT - A component of the set of
dividing points and terminal points of a continuum.
A.A. Mal'tsev
AMS 1980 Subject Classification: 54F15
ADAMS METHOD - A finite-difference method for
solving Cauchy's problem for systems of first-order dif-
ferential equations
y' = I(x,y), y(xo) = Yo
In integration over a grid with a constant step
Xn = X 0 + nh the computational formulas may be based
a) on extrapolation

k step size. However, Adams' method proves to be more
Yn+\ = Yn +h ~:U-A!(Xn-A'Yn-~J;
A=O
or b) on interpolation
k-\
Yn+\ = Yn +h ~ v-A!(xn-A,Yn-A).
A=-\
For a given k, the latter formula gives more accurate
results but a non-linear system of equations must be
solved in order to find Yn + I.
In practical work, an approximation by the former
method is found and one or two corrections are per-
formed according to the formula
k-\
y~tP = Yn +h ~ v-A!(Xn-A,Yn-A)+hvt!(xn+"y~~\),
A=O
which converge pointwise if h I v I I I (3f / 3y) I < l.
The initial conditions Y I, . . . ,Yk for Adams' method,
which are needed to begin the calculations by formula
a), must be found by some special method. The error
in the solution can be written as
x.. H [I. g(x, Ijy(" "(I) dl 1h" , +O(h"'),
where Q(x, t) is the solution of the system
dfJ}:, t) = /y(x,y(xfJ(x, t)
when Q(t, t)=E.
The structure of the term O(h k +2) is such that, for
small values of h, it is usually uniformly small as com-
pared to the main term on large intervals of integra-
tion. This means that Adams' method can be applied
over a large integration interval in the case of an
absolutely-stable solution of the differential problem. In
particular, as distinct from the Milne method, it may be
used for finding stable periodic solutions of differential
equations. The standard Adams procedure for integra-
tion with automatic step selection is much more
involved than the standard Runge - Kutta method, since
the step-changing algorithm is much more involved and
the selection of the initial values of Yk is not standard-
ized.
In the case of the equations y' = -ay, a>O, the
extrapolation formula a) has the form
k
Yn+\ = Yn+h~U-A(-aYn-A).
10.=0
Particular solutions of this equation are Yn = CJln, where
Jl is a root of the equation
If ah '2/ _ I u - A I > 2, then one root of this equation
is I JlI=>\ and the rounding-off errors rapidly
increase. When integrating with automatic step selec-
tion, this results in an unjustified diminution of the
ADDITION
economical than the Runge - Kutta method in most
cases; it was first introduced by J.e. Adams in 1855.
References
[1) BEREZIN, I.S. and ZHIDKOY, N.P.: Computing methods, 2, Per-
gamon Press, 1973 (translated from the Russian).
[2) BAKHVALOY, N.S.: Numerical methods: analysis, algebra, ordi-
nary differential equations, Mir, 1977 (translated from the Rus-
sian).
[3) TnrnONOY, A.N. and GoRBUNOY, A.D.: 'Asymptotic expan-
sions of the error in the difference method of solving Cauchy'S
problem for systems of differential equations', USSR Comput.
Math. Math. Phys. 2, no. 4 (1962), 565-586. (Zh. Vychisl. Mat.
i Mat. Fiz. 2, no. 4 (1962),537-548)
[4) LOZINSKI!, S.M.: 'An estimate of the error of numerical
integration of differential equations 1', Izv. Vyssh. Uchebn.
Zaved. Mat. 5, no. 6 (1958), 52-90 (in Russian).
[5) BAKHYALOY, N.S.: 'On an estimate of the error at numerical
integration of differential equations by Adams' extrapolation
method', Dokl. Akad. Nauk SSSR 104, no. 5 (1955), 683-686
(in Russian).
N. S. Bakhvalov
Editorial comments. The Adams method is a special k
or (k+ 1) multistep method. It yields explicit difference
equations (predicator formula) in the extrapolation case a)
and implicit difference equations (corrector formula) in the
interpolation case b). Certain special case are known as
Adams- Bashforth methods, viz. the series of explicit
methods k =0, Uo = 1 (Euler's: method);
k=1, uo=3/2, U-1 =-1 /2;
k=2, uo=23/12, U-1 =-4/3, u_2=5/12; etc.
A certain series of implicit methods is known as
Adams- Moulton methods, viz. k= 1, V1 = 1/2, Vo =1/2
(the'trapezoidal rule);
k=2, V1 =5/12, vo=2/3, V-1 =-1 /12; etc.
References
[A1) Fox, L. and MAYERS, D.F.: Computing methods for scientists
and engineers, Clarendon Press, 1968.
[A2] ROSSER, J.B.: 'Solving differential equations on a hand-held
programmable calculator', in G.H. Golub (ed.): Studies in
numerical analysis, Math. Assoc. Amer., 1984, pp. 199-242.
AMS 1980 Subject Classification: 65LXX, 65M10,
65N10
ADDmON - One of the basic arithmetic operations.
The result of addition is called the sum. The sum of
two numbers a and b is denoted by a + b, and a and b
are called the summands. Addition of numbers is com-
mutative: a + b = b + a, and assocIative:
(a + b) + c = a + (b + c). The operation inverse to addi-
tion is called subtraction.
Addition is also the name usually given to the opera-
tion in an Abelian group (in additive notation) and to
the binary operation in a ring under which its elements
form an Abelian group. Here also, addition is associa-
tive and commutative. Sometimes, a non-commutative
operation in a group is called addition, such as in the
case of a multi-operator group.
O.A. Ivanova
AMS 1980 Subject Classification: 00A25, 20XX
35
ADDITION OF SETS
ADDITION OF SETS - Vector addition and certain
other (associative and commutative) operations on sets
Ai' The most important case is when the Ai are convex
sets in a Euclidean space Rn.
The vector sum (with coefficients Ai) is defined in a
linear space by the rule
S = ~AiAi = U {~AiXi}'
x,EA j I
where the Ai are real numbers (see [1)). In the space Rn,
the vector sum is called also the Minkowski sum. The
dependence of the volume S on the Ai is connected
with mixed-volume theory. For convex Ai, addition
preserves convexity and reduces to addition of support
functions (cf. Support function), while for C 2 -smooth
strictly-convex Ai eRn, it is characterized by the addi-
tion of the mean values of the radii of curvature at
points with a common normal.
Further examples are addition of sets up to transla-
tion, addition of closed sets (along with closure of the
result, see Convex sets, linear space of; Convex sets,
metric space of), integration of a continual family of
sets, and addition in commutative semi-groups (see [4]).
Firey p-sums are defined in the class of convex bodies
Ai eRn containing zero. When p ;;;;'1, the support func-
tion of the p-sum is defined as C'L./IIf)l / P, where Hi
are the support functions of the summands. For p.;;; - 1
one carries out (-p )-addition of the corresponding
polar bodies and takes the polar of the result (see [2]).
Firey p-sums are continuous with respect to Ai and p.
The projection of a p-sum onto a subspace is the p-sum
of the projections. When p = 1, the p-sum coincides
with the vector sum, when p = -1 it is called the
inverse sum (see [1]), when p = + 00 it gives the convex
hull of the summands, and when p = - 00 it gives their
intersection. For these four values, the p-sum of polyhe-
dra is a polyhedron, and when p = +2, the p-sum of
ellipsoids is an ellipsoid (see [2)).
The Blaschke sum is defined for convex bodies
Ai eRn considered up to translation. It is defined by
the addition of the area functions [3].
The sum along a subspace is defined in a vector space
X which is decomposed into the direct sum of two sub-
spaces Y and Z. The sum of Ai along Y is defined as
zl(z {~(Yz nAi+
where Y z is the translate of Y for which Yz n z = {z }
(see [1 D
References
[I] ROCKAFELLER, R.T.: Convex analySiS, Princeton Univ. Press,
1970.
[2] FlREY, W.J.: 'Some applications of means of convex bodies',
Pacif. 1. Math. 14 (1964), 53-60.
36
[3) FIREY, W.J.: 'Blaschke sums of convex bodies and mixed
bodies', in Proc. Coli. Convexity Copenhagen, 1965,
Copenhagen Univ. Mat. Inst., 1967, pp. 94-101.
[4) DINGHAS, A.: Minkowskische Summen und Integrale. Superad-
ditive Mengenfunktionale. Isoperimetrische Ungleichungen, Paris,
1961.
VP. Fedotov
AMS 1980 Subject Classification: 03EXX, 04AXX,
52A25, 52A20
ADDITION THEOREM for weights - If a Hausdorff
compactum X can be represented as the union over a
set of infinite cardinality ';;;'T of its subspaces of weight
';;;'T, then the weight of X does not exceed 'T. The addi-
tion theorem (which was formulated as a problem in
[I)) was established in [3] for 'T=~o and in [4] in com-
plete generality. Cf. Weight of a topological space.
References
[1) ALEKSANDROV, P.S. and URYSOHN, P.: Memoire sur les espaces
topologiques compacts, Koninkl. Nederl. Akad. Wetensch.,
Amsterdam, 1929.
[2) ENGELKING, R.: General topology, PWN, 1977 (translated from
the Polish).
[3) SMIRNOV, Yu.M.: 'On metrizability of bicompacta, decompos-
able as a sum of sets with a countable base', Fund Math. 43
(1956), 387-393 (in Russian).
[4) ARKliANGEL'sKIi, A.V.: 'An addition theorem for weights of
sets lying in bicompacta', Dokl. Akad Nauk SSSR 126, no. 2
(1959), 239-241 (in Russian).
[5) ARKHANGEL'sKIi, A.V. and PONOMAREV, V.I.: Fundamentals of
general topology: problems and exercises, Reidel, 1984
(translated from the Russian).
A. V Arkhangel'skii
AMS 1980 Subject Classification: 54A25
ADDITIVE ARITHMETIC FUNCflON - An arith-
metic function of one argument that satisfies the fol-
lowing conditions for two relatively prime integers m, n
f(mn) = f(m) +f(n).
An additive arithmetic function is said to be strongly
additive if f(pa)= f(P) for all prime numbers p and all
positive integers a. An additive arithmetic function is
said to be completely additive if the condition
f(mn)= f(m)+ f(n) is satisfied for relatively non-
prime integers m, n as well; ill such a case
f(pa)=af(p)
Examples. The function Q(m), which is the number
of all prime divisors of the number m (multiple divisors
are counted according to their multiplicity), is an addi-
tive arithmetic function; the function w(m), which is
the number of different prime divisors of the number
m, is strongly additive; and the function log m is com-
pletely additive. J.P. Kubilyus
Editorial comments. An arithmetic function is also called
a number-theoretic function.
AMS 1980 Subject Classification: 10A20

ADDITIVE CATEGORY - A category (; in which for
any two objects X and Y an Abelian group structure is
defined on the set of morphisms HOIllQ:(X, y), such
that the composition of morphisms
HOIIlq(X, Y)XHoIIlq(Y, Z) ~ HOIIlq(X, Z)
is a bilinear mapping. Another necessary condition is
that (; includes a null object (zero object, cf. Null
object of a category) as well as the product X X Y of
any two objects X and Y.
In an additive category the direct sum XEB Y of any
two objects exists. It is isomorphic to their product
XX Y. The dual category to an additive category is also
an additive category.
A functor F: (;~(;I from an additive category (; into
an additive category (;1 is said to be additive if, for any
objects X and Y in (;, the mapping
F: HOIIl(;r(X, Y)~HoIllQ:, (F(X), F(Y) is a homomor-
phism of Abelian groups. An additive category is said
to be pre-Abelian if a kernel (cf. Kernel of a morphism
in a category) and a cokernel exist for any morphism.
If for a morphism u: X ~ Y in an additive category
there exists an image Im(u) and a co-image Coim(u),
then there exists a unique morphism
u: Coim(u)~Im(u) such that the morppism u splits as
the composition
X ~ Coim(u) ~ Im(u) ~ Y.
An Abelian category is additive by definition. Exam-
ples of non-Abelian additive categories are the category
of topological modules over a given topological ring
with respect to the morphisms which are continuous
linear mappings, and also the category of Abelian
groups r with a filtration r = r0 ::J r I::J . . . ::J r n = {O}
with respect to the morphisms which are group
homomorphisms preserving the filtration.
References
[1] BUCUR, I. and DELEANU, A.: Introduction to the theory of
categories and functors, Wiley, 1968.
[2] GROTHENDIECK, A.: 'Sur quelques points d'algebrique homolo-
gique', Tohoku Math. J. 9 (1957), 119-121.
[3] GRUSON, L.: 'Completion abelienne', Bull. Sci. Math. (2) 90
(1966), 17-40. 1. V. Dogac
I hev
AMS 1980 Subject Classification: 18A05, 18816,
18E05
ADDITIVE DmSOR PROBLEM - The problem of
finding asymptotic values for sums of the form:
~ Tk, (mfrk,(m +a),]
m~n
(1)
~ Tk,(mfrk,(n -m),
m<n
where Tk(m) is the number of different factorizations of
an integer mink factors, counted according to multi-
plicity. Here kl and k2 are integers ~2, a is a fixed
ADDITIVE GROUP
integer different from zero and n is a sufficiently large
number. In particular 'T2(m)=T(m) is the number of
divisors of the number m. Sums of the form (1) express
the number of solutions of the equations
X I . . . Xk, - YI . . . Yk, = a, (2)
xI . . . Xk, +Y I . . . Yk, = n, (3)
respectively. Particular cases of the additive divisor
problem (k l =k2 =2, kl =2 and k2 =3) are considered
in [1] - [3]. The additive divisor problem with kl =2
and an arbitrary positive integer k2 was solved using
the dispersion method of Yu.V. Linnik [4].
References
[1] INGHAM, A.E.: 'Some asymptotic formulae in the theory of
numbers', J. London Math. Soc. (1) 2 (1927),202-208.
[2] EsTERMAN, T.: 'On the representations of a number as the sum
of two products', Proc. London Math. Soc. (2) 31 (1930), 123-
m.
[3] HooLY, C.: 'An asymptotic formula in the theory of numbers',
Proc. London Math. Soc. (3) 7 (1957), 396-413.
[4] LINNIK, YU.V.: The dispersion method in binary additive prob
lems, Amer. Math. Soc., 1963 (translated from the Russian).
B.M. Bredikhin
Editorial comments. The function T2(m)=T(m) is also
denoted by d(m) or ao(m), cf. [A 1], Sect. 16.7.
References
[A1] HARDY, G.H. and WRIGHT, E.M.: An introduction to the
theory of numbers, Clarendon Press, 1979.
AMS 1980 Subject Classification: 10JXX
ADDITIVE FUNCTION, finitely-additive function (on
sets, on domains) - A real-valued function JL defined on
a system of sets E and such that
I-'< U E;) = ~I-'<E;) (*)
for any finite number of pairwise-disjoint sets
{E I, . . . ,En} of E whose union also belongs to E.
Countably-additive set functions are an important kind
of additive functions (cf. Countably-additive set func-
tion).
A.P. Terekhin
Editorial comments. Suppose that E is a a-algebra on a
set X. Then a non-negative function p. (taking, possibly, the
val ue + 00) on E is an additive (finitely-additive, countably-
additive) measure if it satisfies (*) for an arbitrary (respec-
tively, finite, countable) number of disjoint sets E; in E.
Usually, a measure (sic) is a countably-additive measure.
References
[A1] ROYDEN, H.L.: Real analysis, Macmillan, 1968.
AMS 1980 Subject Classification: 28A 10
ADDITIVE GROUP of a ring - The group formed by
all elements of the given ring with respect to the addi-
tionoperation in the ring. The additive group of a ring
is always Abelian.
o.A. Ivanova
AMS 1980 Subject Classification: 16AXX, 13AXX
37
ADDITIVE NOISE
ADDITIVE NOISE - An interference added to the
signal during its transmission over a communication
channel. More precisely, one says that a given commun-
ication channel is a channel with additive noise if the
transition function Q(y, .) of the channel is given by a
density q(y, y), y E !lJI, YE CfII =!lJI (!lJI and CfII are the
spaces of the values of the signals at the input and out-
put of the channel, respectively) depending only on the
difference y-y, i.e. q(Y,Y)=q(j-y). In this case the
signal ij at the output of the channel can be represented
as the sum of the input signal ." and a random variable
r independent of it, called additive noise, so that
ij=.,,+r.
If one considers channels with discrete or continuous
time over finite or infinite intervals, the notion of a
channel with additive noise is introduced by the rela-
tion ij(t) =.,,(t) +W), where t is in the given interval,
.,,(t), ij(t) and W) are random processes representing
the signals at the input and the output of the channel
with additive noise, respectively; moreover, the process
W) is independent of .,,(t). In particular, if is a nt)
Gaussian random process, then the considered channel
is called a Gaussian channel.
References
[I] R.: Information theory and reliable communication,
GALLAGER,
McGraw-Hill, 1968.
[2] KHARKEVICH, A.A.: Channels with noise, Moscow, 1%5 (in
Russian).
P.L. Dobrushin
v. V. Prelov
Editorial comments. More generally, especially in sys-
tem and control theory and stochastic analysis, the term
additive noise is used for describing the following way noise
enters a stochastic differential equation or observation equa-
tion: dx= f(x, t) dt+ dw, dy= h(x, t) dt+ dv, where wand v
are Wiener noise processes. The general situation of a sto-
chastic differential equation of the form
dx = f(x, t) dt + g(x, t) dw is referred to as having multiplica-
tive noise.
AMS 1980 Subject Classification: 94AXX
ADDITIVE NUMBER THEORY - The branch of
number theory dealing with problems involving the
partitioning (or decomposition) of integers into sum-
mands of a given type, and also with the algebraic and
geometric analogues of such problems concerning alge-
braic number fields and sets of lattice points. These
problems are said to be additive, and the ones usually
considered are partition problems of large integers.
Classical problems in additive number theory
include: the representation of a number as a sum of
four squares, nine cubes, etc. (cf. Waring problem); the
representation of a given number as a sum of not more
than three prime numbers (cf. Goldbach problem); and
the representation of a given number as a sum of one
38
prime number and two squares (cf. Hardy - Littlewood
problem). Problems in additive number theory are
solved by analytic, algebraic, elementary and mixed
methods, and also by methods based on probabilistic
concepts. Depending on the method selected, additive
problems form a part of various branches of number
theory - analytic, algebraic and probabilistic number
theory.
The first systematic results in this field were obtained
in 1748 by L. Euler, who employed power series to
study the partition of integers into positive summands;
in particular, he investigated the problem of partition-
ing a given number into a given number of summands.
Many classical problems in additive number theory
are solved by means of generating functions. This
method was introduced by Euler and forms the basis of
the analytic methods developed by G.H. Hardy, J.E.
Littlewood and LM. Vinogradov. As a first step, assign
to given sequences Ai = {ad (where ai;;;;'O is an integer,
and i = I, 2, ... ) the power series
j;(z) = L zai
O=S;;;;aj<oo
and consider the corresponding generating function
k 00
F(z) = IIf;(z) = Lr(n)z",
i=J "=0
where r(n)=rk,A(n) is the number of representations of
the number n in the form
n = aJ + ... +ako ai E Ai, A = {AIoA2'" .}.
Here r (n) can be calculated using the Cauchy integral.
In Vinogradov's method trigonometric sums
j;(a) = L i7Tiaai,
f F(a)e27Tiaai da,
J
r(n) =
o
are used instead of power series. A main part of r (n ),
located in intervals distributed over neighbourhoods of
rational points, is extracted. Instead of the analytic pro-
perties of F(z) which, in several problems of additive
number theory, necessitate the use of hypotheses
related to the Riemann hypothesis (cf. Riemann
hypotheses), a crucial role in the calculation of r(n) is
played by purely arithmetic estimates of trigonometric
sums by Vinogradov's method, and by the law of the
distribution of prime numbers in arithmetic progres-
sions, obtained by the analytic methods of the theory
of Dirichlet L-functions. It was found that, depending
on k, either r(n)i=O for all n;;;;.l or r(n)i=O for suffi-
ciently large n (n;;;.no(A, or the relation r(n)i=O is
valid for almost all values of n, i.e.
L 1)
n~x
lim
X--,lo'X)
]. ~nlJO
X
= 1,

or, finally, r(n) can be expressed by an asymptotic for-
mula. The smallest number k satisfying one of the pro-
perties above is denoted by g(A), G(A), Go(A), or
ko(A), respectively. If {a;} = {P}, where {P} is the
sequence of prime numbers, and k = 3, then one
obtains Vinogradov's theorem: Every sufficiently large
odd integer can be represented as a sum of three prime
numbers; if k=2, one obtains Chudakov's theorem:
Almost all even integers can be represented as a sum of
two prime numbers.
Certain problems in additive number theory are
solved by studying the structure of the sets obtained by
addition of sequences A; = {a;} given only in terms of
their densities d(A;)=infnA;(n) / n, where
A;(n)=~1
~ai~n
1. If Al = ... =Ak=A, then the posi-
tivity of d(A;) implies that g(A)<oo. The application
of this result to problems in additive number theory in
which sequences of density zero are added, is realized
by constructing a new sequence, with a positive density,
from the given sequence. This is done, mainly, by sieve
methods (cf. Sieve method), by which it can be proved
that d(A;) is positive. It was shown by L.G.
Shnirel'man in this way that positive integers can be
represented as a sum of a bounded number of prime
summands, while Yu.V. Linnik found an elementary
solution of Waring's problem.
The elementary sieve methods of Y. Brun (cf. Brun
sieve) and of A. Selberg (cf. Selberg sieve), if applied to
a number of problems in additive number theory, yield
results which are as yet unobtainable by analytic tools.
However, the most advanced solutions of certain prob-
lems in additive number theory are obtained by com-
bining analytic and elementary methods. In sieve
methods, the principle of sieving the prime numbers
out of the sequence of positive integers (cf. Era-
tosthenes, sieve of) is extended to other sequences as
well. Thus, a simultaneous, suitab~ accurate, sieving of
the prime numbers ,,;;;n 8\ and ";;;n 2, respectively (where
01 < 1 and O2 < 1 are suitably chosen positive constants)
out of the sequences {m} and {2n - m} yields a solu-
tion of the so-called Goldbach - Euler quasi-problem on
the representation of an even integer as a sum of two
numbers one of which has at most k I, while the other
has at most k2 prime factors.
Linnik in 1959 successfully solved the
Hardy- Littlewood problem by the dispersion method
developed by him. He showed [2] that any sufficiently
large integer can be represented as the sum of a prime
number and two squares of integers. The dispersion
method yielded solutions of several so-called binary
problems, viz. how to find the number Q(n) of solutions
of the equation a+/3=n, where 0: and /3 run through
given sequences which are sufficiently well distributed
in arithmetic progressions. Linnik's method involves the
ADDITIVE NUMBER THEORY
use of elementary concepts of probability theory, which
were also employed by P.L. Chebyshev in his proof of
the law of large numbers. It consists of reducing the
given binary equation to a large number of auxiliary
equations for which the expected number of solutions
S;(n) and the true numbers of solutions Q;(n) of the
equations are connected. If it is shown by a computa-
tion of the dispersion that Q;(n) does not differ much
from S;(n) 'in the mean', then Q(n)= ~S;(n) (within
an acceptable error). The dispersion method was also
employed in the study of the general
Hardy- Littlewood equation.
The range of application of the dispersion method
intersects with the range of application of the method
of the large sieve, developed in 1941 by Linnik. The
method makes it possible to sieve out sequences with
the aid of prime numbers, with an increasing number
of discarded residues. Actually, the method of the large
sieve is a consequence of the laws of distribution of
weakly-dependent random variables.
Additive number theory includes problems whose
systematic study belongs to other branches of number
theory: the problem of representing integers by qua-
dratic or higher-degree expressions; and the study of
Diophantine equations, which may be treated in the
framework of general additive number theory.
In modern number theory various branches of addi-
tive number theory are now being intensively
developed. There is a trend to transfer the problems
and methods of additive number theory to arbitrary
algebraic number fields.
References
[1) VINOGRADOV, I.M.: Selected works, Springer, 1985 (translated
from the Russian).
(2) LINNIK, YU.V.: The dispersion method in binary additive prob
lems, Amer. Math. Soc., 1963 (translated from the Russian).
(3) HUA, L.-K.: 'Abschlitzungen von Exponentialsummen und ihre
Anwendung in der Zahlentheorie', in Enzyklopaedie der
Mathematischen Wissenschaften mit Einschluss ihrer Anwen
dungen, Vol. 1, 1959. Helft 13, Teil 1.
(4) OSTMANN, H.H.: Additive Zahlentheorie, Springer, 1956.
(5) CHUDAKOV, N.G.: Uspekhi Mat. Nauk 4 (1938), 14-33.
(6) BREDIKHIN, B.M.: 'The dispersion method and definite binary
additive problems', Russian Math. Surveys 20, no. 2 (1965),
85-125. (Uspekhi Mat. Nauk 20, no. 2 (1965), 89-130)
(7) DAVENPORT, H.: Multiplicative number theory, Springer, 1980.
(8) HALBERSTAM, H. and RICHERT, H.E.: Sieve methods, Acad.
Press, 1974.
B.M. Bredikhin
Editorial comments. Elementary formulas for the number
of partitions of an integer can be found in [A1], Chapt. XIX.
To derive such formulas the method of generating functions
(cf. Generating function) can be used.
The standard reference on the early historical aspects is,
of course, [A2].
References
[A1] HARDy, G.H. and WRIGHT, E.M.: An introduction to the
theory of numbers, Clarendon Press, 1979.
39
ADDITIVE NUMBER THEORY
[A2] DICKSON, L.E.: History of the theory of numbers, 1-3, Chel-
sea, reprint, 1971.
AMS 1980 Subject Classification: 10JXX
ADDITIVE PROBLEMS - Problems in number
theory concerning the decomposition (or partition) of
integers into summands of a given kind. The solution
of classicru additive problems resulted in the develop-
ment of new methods in number theory. The classical
additive problems include:
1) The Goldbach problem on representing an odd
integer, larger than 5, as the sum of three prime
numbers, and the Euler - Goldbach problem on the
representation of an even integer, larger than 2, as the
sum of two prime numbers. These problems were first
posed in 1742.
2) The Waring problem on representing any positive
integer as a sum of s=s(k) non-negative k-th powers,
where k~2 is given. First posed in 1770.
Other additive problems are, for example:
3) The problem on the representation of posItIve
integers as the sum of a bounded number of prime
numbers (the weak Goldbach problem).
4) The Hardy - Littlewood problem on representing
any integer, larger than 1, as the sum of a prime
number and two squares (posed in the nineteen-
twenties).
5) Problems concerning the representation of all suf-
ficiently large numbers as sums of two numbers with a
bounded number of prime factors.
6) Problems concerning the representation of integers
by quadratic forms in three and four variables, and
similar problems.
Additive problems may be solved by analytic, alge-
braic, elementary and mixed methods (cf. Additive
number theory). A large number of additive problems
belong to one of the following two classes:
a) Ternary additive problems of the type
n=a+{3+y, where a and {3 belong to sufficiently
dense sequences of integers which are well distributed
in arithmetic progressions, while y belongs to a
sequence which may be thin, but for which the
corresponding relevant trigonometric sums behave
properly.
b) Binary additive problems of the type n = a +{3,
with a and {3 subject to the same conditions as those
stated in a).
A universal tool for solving ternary additive prob-
lems for sufficiently large n is the general analytic
method of Hardy - Littlewood - Vinogradov: the
method of trigonometric sums (cf. Vinogradov method).
Binary additive problems usually cannot be solved by
such methods, and are solved by different variants of
the elementary sieve method (cf. Sieve method). Espe-
40
cially powerful results are obtained by the method of
the large sieve and by the dispersion method, due to
Yu.V. Linnik. Additive problems of the type 6) above
are also binary. They are studied by special arithmetic-
geometric methods of the theory of quadratic forms.
References
[1] VINOGRADOV, I.M.: The method of trigonometric sums in the
theory of numbers, Interscience, 1954 (translated from the Rus-
sian).
[2] LINNIK, YU.V.: The dispersion method in binary additive prob-
lems, Amer. Math. Soc., 1963 (translated from the Russian).
[3] LINNIK, YU.V.: Ergodic properties of algebraic fields, Springer,
1968 (translated from the Russian).
[4] HUA, L.-K.: 'Abschatzungen von Exponentialsummen und ihre
Anwendung in der Zahlentheorie', in Enzyklopaedie der
Mathematischen Wissenschaften mit Einschluss ihrer Anwen-
dungen, Vol. I, 1959. Heft 13, Teil I.
B.M. Bredikhin
AMS 1980 Subject Classification: 10JXX
ADDITIVE RELATION - A sub module r of the
direct sum A EB B of two modules A and B over some
ring R. An additive relation can thus be regarded also
as a (not necessary single-valued) mapping r: A ~B or,
more exactly, as a 'many-valued' homomorphism, i.e. a
homomorphism rO of the submodule Def r into the
quotient module B /Ind(r) where
Defr = {aEA: 3bEB (a, b)Er},
Ker r = {a EA: (a, O)Er},
Ind r = Ker r -I .
Here, r -\ : B ~A is the relation inverse to r; it consists
of all pairs (b, a)EBEBA such that (a, b)Er. Con-
versely, if a sub module S CA, a quotient module B / L
of the module B and a homomorphism {3: S~B / L are
given, then there also exists a unique additive relation
r : A ~B such that rO = {3.
H two additive relations r:A~B and s:B~C are
given, then, as in the case of other binary relations, it is
possible to define their product, sr: A ~C, which is the
set of all pairs (a, c) EA EB C such that there exists an
element bEB for which (a, b)Er and (b, C)ES. This
multiplication is associative (wherever defined) and,
moreover, the additive relations form a category with
involution r~r -] .
Additive relations are used in natural definitions of
connecting homomorphisms for exact sequences of
complexes. The above considerations are valid not only
in the category of modules, but also in any other
Abelian category.
References
[1] MAC LANE, S.: Homology, Springer, 1963.
[2] PuPPE, D.: 'Korrespondenzen in Abelschen Kategorien', Math.
Ann 148 (1962),1-30.
A. V. Mikhalev
AMS 1980 Subject Classification: 18810

ADDITIVE THEORY OF IDEALS - A branch of
modern algebra. Its principal task is to represent any
ideal of a ring (or of another algebraic system) as the
intersection of a finite number of ideals of special type
(primary, tertiary, primal, uniserial, etc.). The type of
the representation is so chosen that: I) for any ideal
there exists a representation, in other words, some
'existence' theorem holds; 2) the representations chosen
must be unique apart from certain limitations or, in
other words, some 'uniqueness' theorem must hold. The
fundamental principles of the additive theory of ideals
were introduced in the nineteen-twenties and the
nineteen-thirties by E. Noether [I] and W. Krull [2].
All special features of the additive theory of ideals
are clearly manifested in the case of rings. Let R be a
Noetherian ring, i.e. an associative ring with the max
imum condition for ideals. If A is an ideal of R, then
there exists a largest ideal N of R for which N k CA for
some integer k;;;,.l. This ideal N is known as the pri-
mary radical of A (in R) and is denoted by pr(A). An
ideal Q of R is said to be primary if for any two ideals
A and B of R, the condition
AB CQ, A Q ~ B cpr(Q)
is satisfied. The intersection theorem is valid for primary
ideals: The intersection of two primary ideals having
the same primary radical P is itself a primary ideal
with radical P. This theorem is used to prove an
existence theorem: If the ring R is commutative, then
for any ideal A =l=R there exists a representation of A as
the intersection of a finite number of primary ideals Ai:
(1)
such that none of the ideals Ai contains the intersection
of the other ones, and such that the primary radicals
pr(Ai) are pairwise different. Such representations are
known as non-contractible or primarily reduced [1], [4].
The uniqueness theorem holds for such representations:
If (I) holds and
A = B, n ... nBm (2)
is a second primarily-reduced representation of the
ideal A of the ring R, then m =n and pr(Ai)=pr(Bi) for
l,,;;;;i ";;;;n, provided the ideals Bi are suitably renum-
bered.
The additive theory of ideals of Noetherian commu-
tative rings (the classical additive theory of ideals) has
found numerous applications in various branches of
mathematics.
If the ring R is non-commutative, the above-
mentioned existence theorem is no longer valid, but the
uniqueness and intersection theorems still hold. This is
why, ever since the nineteen-thirties, repeated attempts
have been made to find a generalization of classical
primarity to the non-commutative case such that the
ADDITIVE UNIFORM STRUCTURE
existence theorem, too, remains valid. Such a generali-
zation has in fact been found [4], and is known as ter-
tiarity (cf. Tertiary ideal). It was subsequently shown
that, within certain natural limitations, tertiarity is the
only 'good' generalization of the concept of primarity
[6], [7], [8].
During the nineteen-sixties the additive theory of
ideals developed further within the framework of lattice
theory, of systems with fractions and of multiplicative
systems [4], [5], [6]; this stimulated the development of
the additive theory of ideals for non-associative rings,
normal divisors of a group and submodules of a
module.
References
[I) NOETHER, E.: 'Idealtheorie in Ringbereichen', Math. Ann. 83
(1921),24-66.
[2) KRULL, W.: 'Idealtheorie in Ringen ohne Endli-
chkeitsbedingung', Math. Ann. 101 (1929), 729-744.
[3) ZARISKl, O. and SAMUEL, P.: Commutative algebra, I, Springer,
1975.
[4) LESIEUR, L. and CROISOT, R.: Algebre noetherienne noncommu
tative, Gauthier-Villars, 1963.
[5) MURATA, K.: 'Additive ideal theory in multiplicative systems',
J. Inst. Poiytechn. Osaka City Univ. (AJ 10, no. 2 (1959), 91-
115.
[6) ANDRUNAKIEVICH, V.A. and RYABUKHIN, YU.M.: 'The addi-
tive theory of ideals in systems with residuals', Math. USSR
Izv. 1, no. 5 (1967), 1011-1040. (Izv. Akad. Nauk SSSR Ser.
Mat. 31 (1967), 1057-1090)
[7) RILEY, l.A.: 'Axiomatic primary and tertiary decomposition
theory', Trans. A mer. Math. Soc. 105 (1962), 177-201.
[8) GoYAN, I.M. and RYABUKHIN, Yu.M.: 'On the axiomatic addi-
tive theory of Riley ideals', Mat. Issl. 2, no. I (1967), 14-25 (in
Russian).
[9) FUCHS, L.: 'On primal ideals', Proc. A mer. Math. Soc. 1 (1950),
1-6.
[lO) Itogi Nauk. Algebra Topol. Geom. 1965 (1967), 133-180.
V.A. Andrunakievich
Editorial comments. A primary representation is also
called a primary decomposition. More generally one has pri-
mary decompositions for submodules N of modules Mover
a Noetherian ring R. This is a representation
where each Ass(M / OJ) consists of a single prime ideal P"
(By definition Ass(M) for a module M, the set of prime ideals
associated to M, is the collection of all prime ideals P for
which there exists an xEM such that p={rER: rx=O}.)
There is also a corresponding uniqueness theorem, stating
that there is a reduced decomposition, which of course
means that n.J7'='
OJ' C OJ holds for no i E I and that the
Ass(M / OJ)=P, are all different.
References
[A1) BOURBAKI, N.: Aigebre commutative, Hermann, 1961,
Chapt. 3; 4.
AMS 1980 Subject Classification: 13CXX
ADDITIVE UNIFORM STRUCfURE of a topological
skew-field K - The uniform structure of its additive
41
ADDITIVE UNIFORM STRUCTURE
group. A base of neighbourhoods for the uniform struc-
ture of the commuta!ive topological group of K is
formed by the sets V of all pairs (x, y) such that
x - y E V, where V is an arbitrary neighbourhood of
zero. A covering of a topological skew-field K is uni-
form for the additive uniform structure if a covering of
the type {Vy:yEK} refining it can be found, where V
is an arbitrary neighbourhood of zero and
Vy = {x +y: x E V}. In particular, a base of the addi-
tive uniform structure of the real line is formed by all
coverings by intervals of given fixed length. The real
line is the completion of the field of rational numbers
with respect to its additive uniform structure.
The uniform structure of the group Rn , which is
known as its additive uniform structure, is the product
of the uniform structures of its factors R. A base is
constituted by all coverings by (open) spheres of a
given fixed radius.
v. V. Fedorchuk
Editorial comments. A uniform structure is also called a
uniform topology. A base for a uniform structure is also
called a uniformity. See also Uniform space.
References
[A1] ISBELL, l.R.: Uniform spaces, Amer. Math. Soc., 1964.
AMS 1980 Subject Classification: 54E15, 12JXX
ADDITIVITY - The property of magnitudes that can
be defined as follows. The value of a magnitude
corresponding to a whole object is equal to the sum of
the values of the magnitudes corresponding to its parts
for any division of the object into parts. For instance,
additivity of volume means that the volume of a whole
object is equal to the sum of the volumes of its consti-
tuent parts. Cf. Additive function; Countably-additive
set function.
AMS 1980 Subject Classification: 00A05, 00A25
ADELE - An element of the adele group, i.e. of the
restricted topological direct product
IIck.(cO,)
VE V m
of the group Gk with distinguished invariant open sub-
groups Go.. Here Gk is a linear algebraic group, defined
over a global field k, V is the set of valuations (cf.
Valuation) of k, kv is the completion of k with respect
to v E V, and Ov is the ring of integer elements in kv.
The adele group of an algebraic group G is denoted by
GA' Since all groups Gk. are locally compact and since
Go. is compact, GA is a locally compact group.
Examples. 1) If Gk is the additive group k + of the
field k, then GA has a natural ring structure, and is
called the adele ring of k; it is denoted by A k . 2) If Gk
is the multiplicative group k' of the field k, then GA is
42
called the idele group of k (the idele group is the group
of units in the adele ring Ak)' 3) If Gk = GL(n, k) is the
general linear group over k, then GA consists of the ele-
ments g=(gv)E IIVE vGv for which gv EGL(n, Ov) for
almost all valuations v.
The concept of an adele group was first introduced
by C. Chevalley (in the nineteen-thirties) for algebraic
number fields, to meet certain needs of class field
theory. It was generalized twenty years later to alge-
braic groups by M. Kneser and T. Tamagawa [1], [2].
They noted that the principal results on the arithmetic
of quadratic forms over number fields can be con-
veniently reformulated in terms of adele groups.
The image of the diagonal imbedding of Gk in GA is
a discrete subgroup in GA , called the subgroup of princi-
pal adeles. If 00 is the set of all Archimedean valuations
of k, then C = IIc k,
X IIc 0,
A(ao)
VEOO VEOO
is known as the subgroup of integer adeles. If Gk =k',
then the number of different double cosets of the type
GkxGA(oo) of the adele group GA is finite and equal to
the number of ideal classes of k. The naturally arising
problem as to whether the number of such double
classes for an arbitrary algebraic group is finite is con-
nected with the reduction theory for subgroups of prin-
cipal adeles, i.e. with the construction of fundamental
domains for the quotient space GA / Gk . It has been
shown [5] that GA / Gk is compact if and only if the
group G is k-anisotropic (cf. Anisotropic group).
Another problem that has been solved are the cir-
cumstances under which the quotient space GA / Gk
over an algebraic number field has finite volume in the
Haar measure. Since GA is locally compact, such a
measure always exists, and the volume of GA / Gk in
the Haar measure is finite if and only if the group G
has no rational k-characters (cf. Character of a group).
The number T( G) - the volume of GA / Gk - is an
important arithmetical invariant of the algebraic group
G (cf. Tamagawa number). It was shown on the
strength of these results [5] that the decomposition
CA = U CkXjCA(ao)
j=1
is valid for an arbitrary algebraic group G. If k is a
function field, it was also proved that the number of
double classes of this kind for the adele group of the
algebraic group is finite, and an analogue of the reduc-
tion theory was developed [6]. For various arithmetical
applications of adele groups see [4], [7].
References
[I] WElL, A.: Ade/es and algebraiC groups, Princeton Univ. Press,
1961.
[2AJ TAMAGAWA, T.: 'AdeJes', in AlgebraiC groups and discontinuous
subgroups, Proc. Symp. Pure Math., Vol. 9, Amer. Math. Soc.,
1966, pp. 113-121.

[2B] KNESER, M.: 'Strong approximation', in Algebraic groups and
discontinUOWi subgroups, Proc. Symp. Pure Math., Vol. 9, Amer.
Math. Soc., 1966, pp. 187-198.
[3] CASSELS, l.W.S. and FROm.ICH, A. (EDS.): Algebraic number
theory, Acad. Press, 1967.
[4] PLATONOV, V.P.: 'Algebraic groups', J. Soviet Math. 4, no. 5
(1975),463-482. (Itogi Nauk. Algebra Topol. Geom 11 (1973),
5-37)
[5] BoREL, A.: 'Some finiteness properties of adele groups over
number fields', Publ. Math. IHES, no. 16 (1963), 5-30.
[6] HARDER, G.: 'Minkowskische Reduktionstheorie tiber
Funktionenkorpern', Invent. Math. 7 (1969), 33-54.
[7] PLATONOV, V.P.: 'The arithmetic theory of linear algebraic
groups and number theory', Trudy Mat. Inst. Steklov. 132
(1973), 162-168 (in Russian).
[8] WElL, A.: Basic number theory, Springer, 1974.
v.P. Platonov
Editorial comments. Let I be an index set. For each VEl
let Gy be a locally compact group and Oy on open compact
subgroup. The restricted (topological) direct product of the
Gy with respect to the Oy, above denoted by
G = II GvC Oy) ,
YEI
consists (as a set) of all (Xy)EII IGy such that Xy in Oy for
YE
all but finitely many v. The topology on G is defined by tak-
ing as a basis at the identity the open subgroups II
IUy
YE
with Uy an open neighbourhood of Gy for all V and Uy= Oy
for all but finitely many v. This makes G a locally compact
topological group.
AMS 1980 Subject Classification: 12A85, 20GXX,
22E55
ADIABATIC FLOW - A flow for which there is no
heat exchange between the individual parts of the
liquid and the bodies which surround it. If the flow of
a compressed liquid is adiabatic and continuous, then
the entropy of a liquid element is constant. If the
motion is steady, i.e. if the characteristics of the field of
the flow are independent of time, then the sum of the
kinetic energy and the enthalpy for any liquid element
is constant.
References
[l] LANDAU, L.D. and LIFSHITZ, E.M.: Mechanics, Pergamon,
1960 (translated from the Russian). A P 17 k'~
. . ravors
AMS 1980 Subject Classification: 76N10
ADIABATIC INVARIANT - A term of physical origin
whose meaning in mathematics is not unambiguous.
The adiabatic invariant is usually defined as a quantita-
tive characteristic of the motion of a Hamiltonian sys-
tem which remains practically unchanged during an
adiabatic (i.e. very slow in comparison to the time scale
typical of the motion in the system) change of its
parameters. (The change may last for as long as the
values of these parameters - unlike the adiabatic
invariant - undergo significant changes.) Thus, for the
simplest system
ADIABATIC INVARIANT
':i; = v, ~ - -W2(ft)X, (*)
where ( is a small parameter and (,.)(s) is a positive, suf-
ficiently smooth function, the adiabatic invariant is
v
2
I = wx 2 +-.
W
If (")=const, the system (*) describes an ordinary har-
monic oscillator with frequency (,.). Thus, in this case, if
the variation of the parameters of the system is slow in
comparison to the oscillation period, then its energy
(v 2 + ,i x 2 ) / 2 varies proportion'ally to the frequency.
In this example, as well as in most other cases, it is
assumed that if the parameters did not change at all,
the Hamiltonian system under consideration would be
completely integrable and its motion would be quasi-
periodic (in this case simply periodic); other generaliza-
tions are also known. In the publications by many
mathematicians and experts in celestial mechanics on
Hamiltonian systems which are close to being com-
pletely integrable, the term 'adiabatic invariant' is not
used at all, even though the information presented in
these publications seems to indicate that certain magni-
tudes are in certain cases in fact adiabatic invariants.
'Approximate invariance' of an adiabatic invariant
l(t) means that for all t under consideration the differ-
ence I (t) - 1 (0) remains small. (The derivative
dI(t) / dt can be a magnitude of the same order as the
derivatives of the other parameters of the system, but
the variations of the adiabatic invariant are not cum-
mulative with time.) Such approximate invariance may
take place over a very large but finite period of time
(temporal adiabatic invariants) or over the entire infinite
t-axis (stationary or permanent adiabatic invariants [I D.
The notion of 'slow variation of the parameters of the
system' may be rendered more precise in two ways: a)
the Hamiltonian H is an explicit function of the time t
(non-autonomous system), but the derivative dH / dt is
small; b) the system under consideration, with canoni-
cal variables p, q, is a subsystem of a large system with
variables p, q, p' , q', which is itself autonomous and is
/l
such that either p', q' vary slowly, or else their
variation does not materially affect the subsystem.
In view of the above restrictions, the existence of an
adiabatic invariant can be affirmed only under various
additional assumptions, a precise general formulation
of which is difficult [I]. Temporal adiabatic invariants
for systems of the type just described in fact belong to
the field of asymptotic methods of perturbation theory
(if the problem is posed in a more general manner, a
few rigorous results can also be obtained [4]). To prove
the temporal asymptotic invariance of some magnitude
1 (t), the usual procedure is to construct another magni-
tude J(t) such that the values of l(t) oscillate around
J (t), while dJ (t) / dt has a lower order of magnitude
43
ADIABATIC INVARIANT
than the derivatives of the other parameters of the sys-
tem. Thus, in the example (*) for J = I + f.W' xv / w2 (the
prime denotes the derivative of w with respect to the
argument s =f.t) direct differentiation yields
dJ (t) / dt = O(~); this is a proof that I is a temporal
adiabatic invariant. The existence of permanent adia-
batic invariants is open to doubt in case a) above if H
depends on t in some general manner (in particular, in
the absence of periodicity or limits as t ~ + 00). In case
b) above, the problem of permanent adiabatic invari-
ants is connected with small denominators [2].
Historically, adiabatic invariants played an important
role in the quantum theory of Bohr - Sommerfeld, in
which quantizable magnitudes must be adiabatic invari-
ants. This rule lost its importance with the subsequent
advances in quantum mechanics. Adiabatic invariants
are used in modern physics in the study of motion of
charged particles in electromagnetic fields [3]. Here, the
important magnitude is the ratio vi / H, where H is
the magnetic field strength and v ~ is the velocity com-
ponent of the particle lying in a plane perpendicular to
the vector H; this ratio is an adiabatic invariant if the
magnetic field strength does not significantly vary
along the Larmor radius.
It should also be noted that, during an adiabatic
change of state in quantum mechanics, certain magni-
tudes retain their values (e.g. quantum numbers, except
for processes which bring about a degenerate state of
the system). It is therefore possible to speak of adia-
batic invariants in quantum mechanics as well, but this
concept does not play a significant part in quantum
mechanics and is not ordinarily used in that science.
References
[1] MANDEL'SHTAM, L.I., ET AL.: Zh. Russk. Fiz.-Khim. Obshch. 60,
no. 5 (1928), 413-419.
[2] ARNoL'n, V.I.: 'Small denominators and problems of stability
of motion in classical and celestial mechanics', Russian Math.
Survrys 18, no. 6 (1963), 91-191. (Uspekhi Mat. Nauk 18, no. 6
(1963),91-192)
[3] NORTHROP, T.: The adiabatic motion of charged particles, Inter-
science, 1963.
[4] KASUGA, T.: 'On the adiabatic theorem for the Hamiltonian
system of differential equations in the classical mechanics I, II,
III', Proc. Japan Acad. 37, no. 7 (1961), 366-382'D V A
. . nosov
A.P. Favorskir
AMS 1980 Subject Classification: 70HXX, 81 CXX
Amc COHOMOLOGY, 1- - See l-adic cohomology.
AMS 1980 Subject Classification: 14F30
Amc NUMBER, p- - See p-adic number.
AMS 1980 Subject Classification: 12BXX
Amc TOPOLOGY - A linear topology of a ring A in
which the fundamental system of neighbourhoods of
44
zero consists of the powers ~n of some two-sided ideal
~. The topology is then said to be ~-adic, and the ideal
~ is said to be the defining ideal of the topology. The
closure of any set F CA in the ~-adic topology is equal
to nn;;'O(F+~n); in particular, the topology is separ-
able if, and only if, n n;;'
0 ~n = (0). The separable com-
A
pletion A of the ring A in an ~-adic topology is iso-
morphic to the projective limit lim(A / ~n).
<-
The ~-adic topology of an A-module M is defined in
a similar manner: its fundamental system of neighbour-
hoods of zero is given by the submodules ~n M; in the
~-adic topology M becomes a topological A-module.
Let A be a commutat~ve ring with identity with an
~-adic topology and let A be its completion; if ~ is an
i,?eal of finite type, the topology in A is f-adic, and
A
~n = ~n A. If ~ is a maxjmal ideal, then A is a local
ring with maximal ideal ~. A local ring topology is an
adic topology defined by its maximal ideal (an m-adic
topology).
A fundamental tool in the study of adic topologies of
rings is the Artin - Rees lemma: Let A be a commuta-
tive Noetherian ring, let ~ be an ideal in A, let E be an
A-module of finite type, and let F be a submodule of
E. Then there exists a k such that, for any n ~O, the
following equality is valid:
mncmkE nF) = mk+nE nF.
The topological interpretation of the Artin - Rees
lemma shows that the ~-adic topology of F is induced
by the ~-adic topology of E. It follows that the com-
pletion A of a ring A in the ~-adic topology is a flat A
A-module (d. Flat module), that the completion E <;?f
the A-module E of finite type is identical with EAA,
and that Krull's theorem holds: The ~-adic topology of
a Noetherian ring is separable if and only if the set
I + ~ contains no zero divisors. In particular, the
topology is separable if ~ is contained in the (Jacob-
son) radical of the ring.
References
[I] ZARISKl, O. and SAMUEL, P.: Commutative algebra, 2, Springer,
1975.
[2] BOURBAKI, N.: Elements of mathematics. Commutative algebra,
Addison-Wesley, 1972 (translated from the French).
VI. Danilov
AMS 1980 Subject Classification: 54H13, 13JXX
ADJOINT CONNECTIONS - Linear connections r
and r such that for the corre_sponding operators of
covariant differentiation v and v there holds
ZB(X, Y) = B(\lzX, Y)+B(X, \lzY)+2w(Z)B(X, Y),
-
where X, Yand Z are arbitrary vector fields, B(, .) is a
quadratic form (i.e. a symmetric bilinear form), and
wO is a I-form (or covector field). One also says that
 ADJOINT DIFFERENTIAL EQUA nON

V and V are adjoint with respect to B. In coordinate

form (where X, Y, Z~ai' B~bij' W~Wi, v~rt),
akbij - fLbsj - f~As = 2Wkbij'
For the curvature _operators R and R and torsion
- If y(t) and W) are arbitrary solutions of 1(y)=0 and
n~)=O, respectively, then
operators T and T of the connections v and v, n k-l
respectively, the following relations hold: ~ ~(-ly(an_kI)U)y<k-j-l) const, tEl.
k = Ij =0
B(R(U,Z)X, Y)+B(X, R(U, Z)Y) = A knowledge of m (..;,n) linearly independent solutions
= 2{ w([U, Z])- Uw(Z)+Zw(U)}B(X, y), of the equation t(~)=O enables one to reduce the order
B(Z, AT(X, Y) - B(AT(Z, Y)X) = of the equation 1(y)=0 by m (see [1] - [3]).
For a system of differential equations
= B(AT(Z, X), y), AT = T- T.
L(x) = 0, L(x) - x+A(t)x, tEl,
In coordinate form,
where A(t) is a continuous complex-valued (n Xn)-
R':;Am + R~ibjm = -2(a,w, - asw,)bij' matrix, the adjoint system is given by
ATijbsk-ATk)bsi-An,bsj = 0. L'(I/J) - -~+A'(tN = 0, tEl
References (see [l], [4]), where A * (t) is the Hermitian adjoint of
[I] NORDEN, A.P.: Spaces with an affine connection, Nauka,
A(t). The Lagrange identity and the Green formula
Moscow-Leningrad, 1976 (in Russian). Ml 17 - kh
. . ~ o/tse ovsk'-
II take the form
Editorial comments. Instead of the name adjoint con- - -,- d-
nections one also encounters conjugate connections.
(I/J,L(x-(L (I/J), x) = dt(I/J,x),
Sometimes the 1-form w is not mentioned in the notion of
adjoint connections. Strictly speaking this notion of an
'adjoint connection' should be called 'adjoint with respect to
[[(~'L(X-(L'(I/J)'X)ldt = (~,x{~~;
Band w'.
where (', .) is the standard scalar product (the sum of
AMS 1980 Subject Classification: 53805
the products of coordinates with equal indices). If x(t)
ADJOINT DIFFERENTIAL EQUATION to an ordi- and 1/I(t) are arbitrary solutions of the equations
nary linear differential equation 1(y)=0 - The ordinary L(x)=O and L *(1/;)=0, then
linear differential equation n~) = 0, where
(~t), x(t _ const, tEl.
/(y) ao(t)y<n)+ ... + an (t)y, (1)
The concept of an adjoint differential equation is
d>
y< = ~, YEC(I), akEC-k(I), closely connected with the general concept of an adjoint
operator. Thus, if I is a linear differential operator act-
ao(t) =j=. 0, tEl; ing on c n (l) into C(1) in accordance with (1), then its
Crn(I) is the space of m-times continuously- adjoint differential operator t maps the space C* (I)
differentiable complex-valued functions on 1 =(a, {3), adjoint to C(1) into the space C n* (1) adjoint to en(l).
and The restriction of t to en(l) is given by formula (2)
t(~)-(-I)n(ao~)<n)+(-I)n-l(al~)<n-l)+ ... +an~' (2)
(see [5]).
Adjoints are also defined for linear partial differen-
~ E C(l) tial equations (see [6], [5]).
(the bar denotes complex conjugation). It follows at Let ~=[to, tdeI, and let Uk be linearly independent
once that linear functionals on cn(~). Then the boundary value
(/1 +/ 2)' = I; +1;, CAl)' = Xl', problem adjoint to the linear boundary value problem
l(y) = 0, tEA, Uk(y) = 0, k= I, ... ,m, m<2n, (3)
for any scalar A. The adjoint of the equation n~)=O is
1(y)=0. For all n-times continuously-differentiable is defined by the equations
functions yet) and t), Lagrange's identity holds: t'm = 0, Uj(~) = 0, j= I, ... ,2n -m. (4)
-
$(y)-t'(~)y = -
d {~n k-l }
~(-IY(an_kI)u>y<k-j-l) . Here the Uj are linear functionals on cn(~) describing
dt k=lj=O the adjoint boundary conditions, that is, they are defined
It implies Green's formula in such a way that the equation (see Green fonnulas)

1\
T

j[II(y)- I' (~)y I dt = j[II(y)-t(~Yldt =

S 10

45
ADJOINT DIFFERENTIAL EQUATION

holds for any pair of functions y, gEcn(~) that satisfy For the linear boundary value problem
the conditions Uk(y)=O, k=l, ... ,m; Vj(~)=O, L(x) _ x+A(t)x = 0, U(x) = 0, tEt!., (7)
j=l, ... ,2n-m.
If where U is an m-dimensional vector functional on the
Uk(y)
n
= ~[O:kpy(p-I)(to)+,8kpy(p-I)(tl)l space Cn(~) of continuously-differentiable complex-
p=1 valued n-dimensional vector functions with m <2n, the
are linear forms in the variables adjoint boundary value problem is defined by
y(P-l)(to), y(P-I)(t l ), p=I, ... ,n, L*(~) = 0, U*(~) = 0, tEt!. (8)
then Vj mare linear forms in the variables (see [ 1D. Here U* is a (2n - m )-dimensional vector
&-I)(tO), &-I)(tl), p= 1, ... ,no functional defined such that the equation

Examples. For the problem (~t), X(tII:I\

I-to
=0
y+a(t)y = 0, O.,;;;t.,;;;l,
holds for any pair of functions x, "'ECk(~) satisfying
y(O)+ay(1)+,8y(1) = 0, the conditions
y(0)+yy(I)+8y(l) = 0, U(x) = 0, u*(~) = o.
with real a(t), n, {3, y, 8, the adjoint boundary value The problems (7), (8) possess properties analogous to
problem has the form those listed above (see [I D.
~+a(t~ = 0, O";;;t.,;;;l, The concept of an adjoint boundary value problem is
o:O)+yO)+I) = 0, closely connected with that of an adjoint operator [5].
Adjoint boundary value problems are also defined for
,80)+80)+1) = o. linear boundary value problems for partial differential
If problem (3) has k linearly independent solutions equations (see [6], [7]).
(in this case the rank r of the boundary value problem References
is equal to n - k), then problem (4) has m - n + k [1] KAMKE, E.: Differentialgleichungen: LOsungsmethoden und
linearly independent solutions (its rank is LOsungen, 1, Chelsea, reprint, 1971.
r' = 2n - m - k). When m = n, problems (3) and (4) [2] NAiMARK, M.A.: Linear differential operators, Moscow, 1969
(in Russian).
have an equal number of linearly independent solu- [3] CODDINGTON, E.A. and LEVINSON, N.: Theory of ordinary dif-
tions. Therefore, when m = n, problem (3) has only a ferential equatiOns, McGraw-Hill, 1955.
trivial solution if and only if the adjoint boundary [4] HARTMAN, P.: Ordinary differential equations, Birkhauser, 1982.
[5] DUNFORD, N. and ScHWARTZ, J.T.: Linear operators. Spectral
value problem (4) has the same property. The Fredholm theory, 2, Interscience, 1963.
alternative holds: The semi-homogeneous boundary [6] MIKHAiLOv, V.P.: Partial differential equations, Moscow, 1976
value problem (in Russian).
[7] WLADIMIROW, W.S. [V.S. VLADIMIROV]: Gleichungen der
l(y) = f(t), Uk(y) = 0, k=I, ... ,n, mathematischen PhYSik, Deutsch. Verlag Wissenschaft., 1973
(translated from the Russian).
has a solution if J(t) is orthogonal to all non-trivial E.L. Tonkov
solutions t) of the adjoint boundary value problem AMS 1980 Subject Classification: 34A30, 34805
(4), i.e. if 1\

j"&t)f(t) dt = 0 ADJOINT FUNCfOR - A concept expressing the

10 universality and naturalness of many important
(see [1] - [3], [7]). mathematical constructions, such as a free universal
For the eigen value problem algebra, various completions, and direct and inverse
limits.
l(y) = '>"y, Uk(y) = 0, k=l, ... ,n, (5)
Let F: S1'~G' be a covariant functor in one argument
the adjoint eigen value problem is defined as from a category S1' into a category CL F induces a func-
I'm =~, Vjm = 0, J=l, ... ,n. (6) tor
If A is an eigen value of (5), then Jl=X is an eigen value
of (6). The eigen functions yet), t)
corresponding to
where ~* is the category dual to st', is the category of
sets, and H dX, Y): st'* X G'~ is the basic set-valued
eigen values A, Jl of (5), (6), respectively, are orthogonal
functor. The functor HF is contravariant in the first
if A=FJl (see [l] - [3]): argument and covariant in the second. Similarly, any
1\
covariant functor G: G'~~ induces a functor
jy(t)t)dt = 0.
to H c;CX, Y) = H\l(X, G( Y: sf X IT __ ;S,

46

which is also contravariant in the first argument and
covariant in the second. The functors F and G are
adjoint, or form an adjoint pair, if HF and HG are iso-
morphic, that is, if there is a natural transformation
(): HF ~HG that establishes a one-to-one correspon-
dence between the sets of morphisms H \i;(F(X), Y) and
H Sf (X, G(Y)) for all objects X EOb Sf and Y EOb (.
The transformation () is called the adjunction of F with
G, F is called the left adjoint of G and G the right
adjoint of F (this is written (): FG, or simply FG). The
transformation ()-I: HG~HF IS called the co-
adjunction.
Let ():FG. For all XEObSf and YEOb(, let
x = O(IF(X), 1jy = 0- 1(lG(y'
The families of morphisms {(x} and {1jy} define
natural transformations (: IdSf~GF and 1j: FG~Ido:,
called the unit and co-unit of the adjunction (). They
satisfy the following equations:
G(1jY}G(Y) = IG(y), 1jF(x)F(x) = IF(X)'
In general, a pair of natural transformations
cp: IdSf~GF and 1/;: FG~Id\i; leads to an adjoint pair (or
adjunction) if the following equations hold:
G(I/IY)</>G(Y) = IG(Y)' I/IF(x)F(c[>x) = IF(X)
for all objects X and Y. A natural transformation
cp: IdSf~GF is the unit of some adjunction if and only
if for any morphism a: X~G(Y) in Sf there is a unique
morphism a':F(X)~Yin ( such that a=G(a')x. This
property expresses the fact that F(X) is a free object
over X with respect to the functor G in the sense of the
following definition. An object Y E Ob ( together with
a morphism (: X~G(Y) is free over an objectXEOb Sf if
every morphism a: X~G(Y') can be written uniquely
in the form a = G( a')( for some morphism a' : Y ~ y'. A
functor G: (~Sf has a left adjoint if and only if for
every X EOb Sf there is an object Y that is free over X
with respect to G.
Examples of adjoint functors. 1) If G: (~, where
is the category of sets, then G has a left adjoint if and
only if it is representable. A representable functor
Gc:::::.HA =H dA, Y) has a left adjoint if and only if all
co-products IIxExAx exist in (, where X EOb and
Ax =A for all x EX.
2) In the category of sets, for any set A the basic
functor HA(Y)=H(A, Y) is the right adjoint of the
functor XXA.
3) In the category of Abelian groups, the functor
Hom(A, Y) is the right adjoint of the functor XA of
tensor multiplication by A, and the imbedding functor
of the full subcategory of torsion groups is the left
adjoint of the functor of taking the torsion part of any
Abelian group.
4) Let P: ~~ be the forgetful functor from an
ADJOINT GROUP
arbitrary variety of universal algebras into the category
of sets. The functor P has a left adjoint F: ~~,
which assigns to every set X the free algebra of the
variety ~ with X as set of free generators.
5) The imbedding functor Ido:, Sf: (~Sf of an arbi-
trary reflexive subcategory ( of a category Sf is the
right adjoint of the (-reflector. In particular, the
imbedding functor of the category of Abelian groups in
the category of groups has a left adjoint, which assigns
to every group G its quotient group by the commutator
subgroup.
Properties of adjoint functors. The left adjoint functor
of a given functor is uniquely determined up to isomor-
phism of functors. Left adjoints commute with co-
limits (e.g. co-products) and send null objects and null
morphism into null objects and null morphisms, respec-
tively.
Let Sf and ( be categories that are complete on the
left and locally small on the left. A functor G: @~Sf
has a left adjoint F: Sf~( if and only if the following
conditions hold: a) G commutes with limits; b) for
every X EOb Sf, at least one of the sets H(X, G(y)),
Y EOb Sf, is non-empty; and c) for every X EOb Sf,
there is a set S cOb ( such that every morphism
a: X ~G(Y) is representable in the form a = G(a')cp,
where cp: X ~G(B), BE S, a' : B ~ Y.
By passing to dual categories, one may establish a
duality between the concepts of a 'left adjoint functor'
and a 'right adjoint functor'; this enables one to deduce
the properties of right adjoints from those of left
adjoints.
The concept of an adjoint functor is directly con-
nected with the concept of a triple (or monad) ill a
category.
References
(I] TSALENKO, M.SH. and SHUL'GE1FER, E.G.: Fundamentals of
category theory, Moscow, 1974 (in Russian).
[2] MACLANE, S.: Categories for the working mathematician,
Springer, 1971. M.Sh. Tsalenko
Editorial comments. A category is called complete on
the left if small diagrams have limits. A category is called
locally small on the left if it has small hom-sets. The state-
ment that a functor has a left adjoint if and only if a), b) and
c) above holds, is called the Freyd adjoint functor theorem.
AMS 1980 Subject Classification: 18A10
ADJOINT GROUP of a group G - The linear group
Ad G that is the image of the Lie group or algebraic
group G under the adjoint representation (d. Adjoint
representation of a Lie group). The adjoint group Ad G
is contained in the group Aut g of automorphisms of
the Lie algebra g of G, and its Lie algebra coincides
with the adjoint algebra ad g of g. A connected semi-
simple group is a group of adjOint type (i.e. is iso-
47
ADJOINT GROUP
morphic to its adjoint group) if and only if its roots
generate the lattice of rational characters of the maxi-
mal torus; the centre of such a group is trivial. If the
ground field has characteristic zero and G is connected,
then Ad G is uniquely determined by the Lie algebra g
and is either called the adjoint group or the group of
inner automorphisms of g. In particular, if G i1: serni-
simple, Ad G coincides with the connected component
of the identity in Aut g.
References
[1) PONTRYAGIN, L.S.: TopolOgical groups, Princeton Univ. Press,
1958 (translated from the Russian).
(2) SERRE, J.-P.: Lie algebras and Lie groups, Benjamin, 1965
(translated from the French).
(3) HUMPHREYS, J.E.: Linear algebraic groups, Springer, 1975.
A.L. Onishchik
Editorial comments.
References
[A1] BOURBAKI, N.: Elements of mathematics. Lie groups and Lie
algebras, Addison-Wesley, 1975, Chapt. 2; 3 (translated
from the French).
AMS 1980 Subject Classification: 20G05, 22E10,
22E15,22E20
ADJOINT LINEAR TRANSFORMATION of a linear
transformation A - The linear transformation A' on a
Euclidean space (or unitary space) L such that for all
x, y EL, the equality
(Ax, y) = (x, A 'y)
between the scalar products holds. This is a special case
of the concept of an adjoint linear mapping. The
transformation A' is defined uniquely by A. If L is
finite-dimensional, then every A has an adjoint A', the
matrix fJB of which in a basis e 1, . . . ,en is related to
the matrix .;;I of A in the same basis as follows:
fJ6 = G~I.9l'G,
where .;;I' is the matrix adjoint to .;;I and G is the
Gram matrix of the basis e 1, . . . ,en.
In a Euclidean space, A and A' have the same
characteristic polynomial, determinant, trace, and eigen
values. In a unitary space, their characteristic polyno-
mials, determinants, traces, and eigen values are related
by complex conjugation.
TS. Pigolkina
Editorial comments. More generally, the phrase 'adjoint
transformation' or 'adjoint linear mapping' is also used to
signify the dual linear mapping cf,": M' --,>L ' of a linear map-
ping cp: L--,>M. Here M' is the space of (continuous) linear
functionals on M and cp'Cm'CI))=m'CCP(/)). The imbeddings
L--,>L', M--,>M', 1...... (-./) connect the two notions. Cf. also
Adjoint operator.
References
[A 1] REED, M. and SIMON, B.: Functional analysis, 1, Acad. Press,
1972, Sect. 2.
AMS 1980 Subject Classification: 15A04, 47 A05
48
ADJOINT MATRIX, Hermitian adjoint matrix, of a
given (rectangular or square) matrix A = II aik II over
the field C of complex numbers - The matrix A' whose
entries a;k are the complex conjugates of the entries aki
of A, i.e. a;k =aki. Thus, the adjoint matrix coincides
with its complex-conjugate transpose: A' = (A ') where-
denotes complex conjugation and the ' denotes transpo-
sition.
Properties of adjoint matrices are:
(A+B)" = A'+B', (M)' = M',
(AB)' = B'A', (A')~l = (A~l)", (A')' = A.
Adjoint matrices correspond to adjoint linear
transformations of unitary spaces with respect to ortho-
normal bases.
For references, see Matrix.
TS. Pogolkina
AMS 1980 Subject Classification: 15A04, 10C05
ADJOINT MODULE, contragradient module, dual
module - The module of homomorphisms of a given
module into the ground ring. More precisely, let M be
a left module over a ring R. The Abelian group
HomR(M, R) of homomorphisms of Minto R regarded
as a left R-module can be made into a right R-module
M' by putting
x(cpA) = (x </A, xEM, </>EHomR(M, R), AER.
This right module M' is called the adjoint of M. For
x EM, one can define an element x EM" by putting
x(cp)=x(cp) for all cpEM'. This defines a homomor-
phism of Minto M'. For any left R-module C, the
mapping~: M' RC~HomR(M, C) given by
x((cpc)n = (x </c, xEM, </>EM', CEC,
is also a homomorphism. Both of these are isomor-
phisms when M is a finitely-generated projective
module [2]. It follows from the properties of the functor
S II
Hom that (~M ~ M" (where ~ is the direct
sum, and II
the direct product) and that there is a
homorphism of M" into M'. The composite map-
ping M' ~M'" ~M' is the identity, but M" need
not be isomorphic to M. The torsionjree modules in
the sense of Bass are those for which the above
homomorphism of Minto M" turns out to be a
monomorphism. This property is equivalent to the
imbeddability of M in a direct product of copies of the
ground ring. If R is right and left Noetherian, then the
mapping Mf->M' defines a duality between the
categories of finitely-generated left and right R-modules
if and only if R is a Quasi-Frobenius ring.
References
[1] BOURBAKI, N.: Elements of mathematics. Algebra: Modules.
Rings. Forms, 2, Addison-Wesley, 1975, Chapt. 4; 5; 6
(translated from the French).

[2] MAc LANE, S.: Homology, Springer, 1963.
[3] MISHINA, A.P. and SKORNYAKOV, L.A.: Abelian groups and
modules, Arner. Math. Soc., 1976 (translated from the Russian).
L.A. Skornyakov
AMS 1980 Subject Classification: 13C99, 16A53
ADJOINT OPERATOR A linear operator
A': Y' ~X' (where X' and y' are the strong duals of
locally convex spaces X and Y, respectively), con-
structed from a linear operator A : X ~ Y in the follow-
ing way. Let the domain of definition DA of A be
everywhere dense in X. If for all XEDA,
<Ax, g> = <x, g'>,
where AXEY, gEY' and g'EX', then A'g=g' is a
uniquely defined operator from the set DA' of elements
g satisfying (*) into X'. If DA = X and A is continuous,
then A' is also continuous. If, in addition, X and Y are
normed linear spaces, then II A' II = II A II. If A is com-
pletely continuous, then so is A'. Adjoint operators are
of particular interest in the case when X and Yare Hil-
bert spaces.
References
[I] YOSIDA, K.: Functional analysis, Springer, 1980.
[2] RIEsz, F. and SZOKEFALVI-NAGY, B.: Lefons d'analyse fonc-
tionelle, Acad. Sci. Hongrie, 1953.
v.I. Sobolev
Editorial comments. In Western literature the adjoint
operator as defined above is usually called the dual or con-
jugate operator. The term adjoint operator is reserved for
Hilbert spaces, in which case it is defined by
(Ax, g) = (x, A' g),
where (-, .) denotes the Hilbert space inner product.
References
[A1] TAYLOR, A.E. and LAY, D.C.: Introduction to functional
analysis, Wiley, 1980.
AMS 1980 Subject Classification: 47 A05
ADJOINT REPRESENTATION OF A LIE GROUP
or algebraic group G - The linear representation Ad of
G in the tangent space Te( G) (or in the Lie algebra g of
G) mapping each a E G to the differential
Ada = d(Inta)e of the inner automorphism
Int a : x ~axa -I. If G <: GL( V) is a linear group m a
space V, then
(Ada)X = aXa- 1 , XETe(G) =g C End(V).
The kernel Ker Ad contains the centre of G, and if G is
connected and if the ground field has characteristic
zero, coincides with this centre. The differential of the
adjoint representation of G at e coincides with the
adjoint representation ad of g.
The adjoint representation of a Lie algebra g is the
linear representation ad of the algebra g into the
module g acting by the formula
(adx)y = [x,y], X,yEg,
ADJOINT SURFACE
where [ , ] is the bracket operation in the algebra g.
The kernel Ker ad is the centre of the Lie algebra g.
The operators ad x are derivations of g and are called
inner derivations. The image ad g is called the adjoint
linear Lie algebra and is an ideal in the Lie algebra
Der g of all derivations of g, moreover Der g / ad g is
the one-dimensional cohomology space HI(g, g) of g,
defined by the adjoint representation. In particular,
ad g = Der g if g is a semi-simple Lie algebra over a
field of characteristic zero.
References
[I] JACOBSON, N.: Lie algebras, Interscience, 1962.
(*) [2] PONTRYAGIN, L.S.: Topological groups, Princeton Univ. Press,
1958 (translated from the Russian).
[3] SERRE, J.-P.: Lie algebras and Lie groups, Benjamin, 1965
(translated from the French).
[4] HUMPHREYS, J.E.: Linear algebraic groups, Springer, 1975.
A.L. Onishchik
Editorial comments.
References
[A1] BOURBAKI, N.: Elements of mathematics. Lie groups and Lie
algebras, Addison-Wesley, 1975, Chap!. 2; 3 (translated
from the French).
AMS 1980 Subject Classification: 17805, 20G05,
22E10, 22E15, 22E20
ADJOINT SPACE of a topological vector space E -
The vector space E' consisting of continuous linear
functions on E. If E is a locally convex space, then the
functionals fEE' separate the points of E (the
Hahn- Banach theorem). If E is a normed space, then
E' is a Banach space with respect to the norm
Ilfll = sup I [ex) I.
x#! II x II
There are two (usually different) natural topologies on
E' which are often used: the strong topology deter-
mined by this norm and the weak-*-topology.
References
[I] RAiKov, D.A.: Vector spaces, Noordhoff, 1965 (translated from
the Russian).
Editorial comments. Instead of the term adjoint space
one more often uses the term dual space. The weak-*-
topology on E' is the weakest topology on E' for which all
the evaluation mappings f..... f(z), fEE', XEE, are continu-
ous.
References
[A 1] ScHAEFER, H.H.: Topological vector spaces, Macmillan,
1966.
AMS 1980 Subject Classification: 46A20, 46810
ADJOINT SURFACE - A surface Y that is in Peter-
son correspondence with a given surface X and is,
moreover, such that the asymptotic net on Y
corresponds to a conjugate net a on X with equal
invariants, and vice versa. The adjoint surface Y is the
49
ADJOINT SURFACE
rotation indicatrix for X, and vice versa. If (J is a princi-
pal base for a deformation of X, then Y is a Bianchi
surface.
I.Kh. Sabitov
AMS 1980 Subject Classification: 53A05
ADJUSTMENT MEfHOD - A method in which the
solution u of a stationary problem
Au =j
is regarded as the steady-state limit solution for t~oo
of a Cauchy initial value problem for a non-stationary
evolution equation involving the same operator A (cf.
Cauchy problem). This evolution equation may e.g. be
of the form m diu(t)
i~l Ci dt i = j - Au(t),
dku
dt k I = UOb k=O, ... ,m-I.
1=0
Here the Ci are suitable operators which guarantee the
existence of the 'adjustment limit' liml-+OO u(t)=u.
A result of using adjustment is that it permits one to
use approximate solution methods of (2) in order to
construct iteration algorithms for solving equation (1)
(cf. Iteration algorithm). Thus, for the non-stationary
equation (2) one could employ a discretization (dif-
ferencing) with respect to t solution method which is
convergent and stable to obtain approximate solutions.
For example, for m = 1, an explicit method of the form
C 1 U(tn+l;-U(t n) = j-Au(tn)
n the pressure p, after which they leave the surface (are
where Tn =tn +3 - tn >0. And then this method can be
interpreted as an iteration algorithm
C1(un+1-u n) = Tn(j-Au n), n =0,1, ... , UO = uoo,
for solving equation (1), in which C I and Tn are now
seen as characterizing this (iteration) method.
Varying the form of the operators Ci and considering
different discretizations with respect to t in equation (2)
(explicit schemes, implicit schemes, splitting schemes,
etc.) gives the possibility of obtaining a wide variety of
iteration methods for solving equation (1). For these
methods equation (2) will be the closure of the compu-
tational algorithm (d. Oosure of a computational algo-
rithm). A generalization of the adjustment method IS
the continuation method (to a parametrized family).
References
[1] BAKHVALOY, N.S.: Numerical methods: analysis, algebra, ordi
nary differential equations, Mir, 1977 (translated from the Rus
sian).
[2] GoDUNOY, S.K. and RYABEN'Kli, V.S.: The theory of difference
schemes, North-Holland, 1964.
[3] MARCHUK, G.!. and LEBEDEY, V.I.: Nt:merical methods in the
theory of neutron transport, Moscow, 1971 (in Russian).
VI. Lebedev
50
Editorial comments. The adjustment method is also
called the time-stepping method. Such a time-stepping
method with m =1 arises quite naturally when an elliptic
boundary value problems is viewed as a steady state of a
(dissipative) parabolic problem (ct. [A2]). Because of the
inherent (numerical) stiffness an implicit discretization
method, such as BDF, should be advocated.
Similar ideas can be used in optimization theory by e.g.
constructing a mechanical system whose stable equilibrium
(1) state is the desired optimum. Again this leads to all kinds of
iteration algorithms, as in [A3].
References
[A 1] BABUSHKA, I. and SoBOLEY, S.L.: 'The optimization of numeri-
cal processes', Appl. Mat. 10 (1965),96-130.
[A2] KUBICEK, M. and HLAVACEK, Y.: Numerical solution of non-
linear boundary value problems with applications, Prentice-
(2) Hall,1983.
[A3] RAzUMIKHIN, B.S.: Physical models and equilibrium methods
in programming and economics, Reidel, 1984.
[A4] RHEINBOLDT, W.e.: Numerical analysis of parametrized non-
linear equations, Wiley, 1986.
AMS 1980 Subject Classification: 65HXX
ADSORPTION - 'Consumption' of matter from a gas
or from an opening of the interface between them (or
from the surface of a solid body). In other words,
adsorption is the 'consumption' by an adsorbate from a
volume of gas on the surface of the adsorbent. Adsorp-
tion is a particular case of sorption.
The molecules of the adsorbate falling on the surface
of the adsorbent are retained by the surface force for a
period of time, which depends on the natures of the
adsorbent and adsorbate, on the temperature T and on
desorbed). Under conditions of thermodynamic and
molecular equilibrium, the rates of adsorption and
desorption are equal. The relation between the relative
pressure cf>=p / ps of the adsorbent and the relative
concentration O=c / cs , where the index s stands for
the limit value at a constant temperature, is known as
the adsorption isotherm.
Langmuir's equation of mono-molecular adsorption
has the form
8
</> = k(l-8)'
where k is the equilibrium constant which roughly
describes the interaction between the adsorbent and the
adsorbate.
Brunauer's equation [l] is commonly used in the case
of a homogeneous surface of the adsorbent and poly-
molecular adsorption.
Posnov's formula [2], which is empirical, is widely
employed for capillary bodies:
7i1 = A In</>+ 1,
where A is a coefficient which varies with the tempera-
ture and with the structure of the adsorbent.

References
[I) BRUNAUER, S.: Adsorption of gases and vapors, Princeton Univ.
Press, 1943.
[2) POSNOV, V.A.: Zh. Tekhn. Fiz., no. 23 (1953), 865.
[3) IL'IN, B.V.: The nature of adsorption forces, Moscow-Leningrad,
1952 (in Russian).
[4) BOER, J.H. DE: The dynamical character of adsorption, Claren-
don Press, 1968.
A. V. Lykov
Editorial comments. The Brunauer- Emmett- Teller
equation, or BET-equation, [A1] generalizes the Langmuir
equation. It assumes multi-molecular adsorption for which
each layer obeys a Langmuir equation. There are several
modifications, cf. [A2].
References
[A1] BRUNAUER, S., EMMETT, P.H. and TELLER, E.: J. Amer. Chem.
Soc. 60 (1938), 309.
[A2] BRUNAUER, S. and COPELAND, L.E.: 'Surface tension, adsorp-
tion', in E.U. Condon and H. Odishaw (eds.): Handbook of
physics, Vol. 2, McGraw-Hili, 1967, Chapt. 7.
AMS 1980 Subject Classification: 76S05, 80A20,
82A99
AERODYNAMICS, MATHEMATICAL PROBLEMS
OF - Problems involved in the solution of fundamental
equations in aerodynamics, which are an accurate
description of the laws of motion of a gas-like medium
and its interaction forces with solids moving in this
medium. Turbulence, for which no satisfactory
mathematical model has yet been established, forms an
exception to this rule. Processes without turbulence or
in which the turbulence can be neglected may, in prin-
ciple, be sufficiently completely investigated by a
mathematical solution of the relevant equations. The
real possibilities of solving such equations differ from
the theoretical; accordingly, the orders of the individual
terms in the equations are estimated, and on this basis
simplified mathematical models are established for dif-
ferent classes of aerodynamic processes, in which one
or more 'small' terms in the complete set of aero-
dynamic equations are discarded. Even for such simpli-
fied models analytic solutions can be obtained only for
the simplest cases, and numerical methods are widely
employed in practical work in aerodynamics. These
methods have become prevalent with the introduction
of electronic computers.
As regards the nature of the equations involved in
aerodynamical problems, these can be subdivided into
four large classes.
1) Problems in the aerodynamics of an ideal (non-
viscous and thermally non-conducting) fluid, the word
'fluid' being used in this context to denote both liquids
and gases.
2) Problems in the aerodynamics of viscous fluids.
3) Problems in the aerodynamics of radiant gases.
4) Problems in the aerodynamics of rarefied gases.
The equations of an ideal fluid are obtained from the
AERODYNAMICS. MATHEMATICAL PROBLEMS OF
general equations by making the Reynolds number -
which describes the order of the ratio between the iner-
tial and the viscous forces - grow without limit.
In most problems of applied aerodynamics of aircraft
and rockets, and in seacraft hydrodynamics, the Rey-
nolds numbers are very large (l(f - 108 ), so that the
model of an ideal fluid yields an accurate description of
flow processes around solid bodies under these condi-
tions, except for the boundary layer zone. The system
of aerodynamic equations of an ideal fluid are:
p-
dV = pg-Vp
dt
(equation of motion);
!!e.
dt
+ p div V = 0
(continuity equation);
p![';2+V) = p(g, V)-div(pV)
(equation of conservation of energy). Here V is the
velocity vector, p is the density, p is the pressure, g is
the acceleration due to gravity and U is the internal
energy. It is a system of equations of the first order
with respect to the functions sought. This is also true
for a mixture of chemically interacting gases: The equa-
tions of the transfer of the mixture components are
added to the above equations, while the heat generation
due to chemical reactions is added to the energy equa-
tion.
A number of simplified models can be derived from
the model of an ideal fluid. The most important ones
are the model of a weight-less fluid and the model of
an incompressible fluid. The former model is the limit-
ing case of the general model as the Froude number
increases without limit; the Froude number describes
the order of the ratio of the inertial to the gravity
forces. This model is very accurate when applied to the
aerodynamics of aircraft and rockets (with Froude
numbers of order 1()2 - 103 ). However, the force of
gravity becomes significant in seacraft hydrodynamics
and especially so in meteorological applications (with
Froude number of order ~ 1).
The equation of an incompressible fluid is obtained
from the general system by putting p = const. This form
of the equation of motion is obtained as the Mach
number tends to zero (which is the relationship of the
characteristic velocity of sound in a fluid medium). In
the general case this is not true; it is only valid for
Froude numbers ~ 1 (i.e. if MaFr), and also for
quasi-stationary processes, in which the Strouhal
number (the criterion of similitude of non-stationary
motions of the liquid or the gas) is of order ~ 1.
The simplest model is that of an incompressible
51
AERODYNAMICS, MATHEMATICAL PROBLEMS OF
fluid. In this model the entire system of equations is
reduced to the single Laplace equation for the velocity
potential cI> (\7 cI> = V). However, it is only in the case
of plane-parallel motions of a weight-less fluid and
axi-symmetric motion around solids of revolution that
the solution of aerodynamic problems is reduced to
problems in classical potential theory.
The problem of flow around a given shape is solved
by classical methods of the theory of conformal map-
ping. The flow around solids of revolution at zero angle
of incidence is determined by the classical solution of
the exterior Neumann problem for the Laplace equa-
tion. In view of the practical importance of these prob-
lems (determination of the characteristics of wing pro-
files and fuselages) many methods for their solutions
have been developed that achieve an acceptable accu-
racy by means of relatively simple calculations. If elec-
tronic computers are employed, there is no point in
using these methods, since accurate numerical methods
for solving these problems belong to the simplest prob-
lems of computational mathematics.
In the spatial case the formal solution of the Neu-
mann problem for the velocity potential will
correspond to the physical picture of the flow only in
very exceptional cases. In a viscous fluid a solid is fol-
lowed by a turbulent trail. As the Reynolds number
increases, this trail becomes thinner (for a continuous
flow) and in the limit it becomes an infinitely thin tur-
bulent surface, the intensity of the turbulent layer of
which vanishes only in exceptional cases (e.g. in the
case of plane-parallel motion). Accordingly, the prob-
lems involving discontinuous velocity potentials in the
region of the flow zone represent the real spatial prob-
lems of an incompressible fluid. The location of this
discontinuity surface is unknown, and for this reason
the exact problem of flow around solid bodies in the
presence of a discontinuous potential surface is a very
complicated non-linear problem. It is only in the linear
approximation, i.e. on the assumption that the solid
around which the fluid is flowing disturbs the basic
mainstream only to a small extent (a thin wing at a
small angle of incidence) that the problem could be
solved. In such a case the surface of turbulence can be
considered as horizontal, and it may be assumed that
the potential along the stream has a constant discon-
tinuity. The normal derivatives of the velocity poten-
tial are given on the projection plane of the solid. Ana-
lytic solutions have been obtained on these asumptions
(for a circular and an elliptic wing). Numerical methods
of solving this problem for other shapes of wings are
available and simplify the task to a considerable extent
(cf. Wing theory).
The principal difficulty in problems involving a
heavy fluid is that the boundary conditions for the
potential on the free surface are non-linear. Exact
52
solutions have been found for planar problems of the
theory of waves only. Little progress has been made on
the spatial problem of waves with a finite amplitude.
On the other hand, the linearized problem (small wave
amplitudes) has been largely solved, and the calculation
of the wave resistance to the motion of a ship is actu-
ally based on the linear theory.
The possibility of aircraft travelling at near-sonic and
supersonic speeds stimulated the search for solutions of
aerodynamic problems of compressible fluids. The
problems involving subsonic, transsonic and supersonic
velocities are all treated in substantially different
manners. In the subsonic range the aerodynamic equa-
tions remain elliptic and the solutions qualitatively
resemble those obtained for incompressible fluids.
In the first period of research on the aerodynamics of
compressible fluids much research was done on the
effect of compressibility on the aerodynamic charac-
teristics of the solids around which the flow takes place.
All these methods involve some sort of linearization of
the velocity of the unperturbed flow. Khristianovich's
method is an exception, permitting as it does to obtain
accurate solutions of the equations of subsonic aero-
dynamics, though for a solid whose shape is unknown
in advance.
In practice, the difficulties involved in the problems
of subsonic aerodynamics are the same as those to be
dealt with when investigating incompressible fluids:
problems involving plane-parallel and axi-symmetric
flows may be fairly easily solved by numerical methods
(cf. Integral-relation method; Adjustment method). In
the case of spatial flows with discontinuous velocity
potential, solutions can only be obtained for linearized
equations. These equations can be converted, by an ele-
mentary transformation, into equations for incompres-
sible fluids in which the boundary conditions and the
conditions on the surface of discontinuity are the same
as those for incompressible fluids.
The zone of trans sonic aerodynamics is, from the
point of view of mathematics, including numerical
mathematics, the most difficult to study. The presence
of local supersonic (parabolic) zones which almost
always terminate in shock waves precludes any analytic
approximation. Moreover, these problems display the
principal drawback of equations of elliptic type - the
propagation of the effect of any particular disturbance
throughout the space. The most suitable method for
solving the problems of transsonic aerodynamics is the
adjustment method, in which a non-stationary aero-
dynamic problem is solved (starting from an arbitrary
initial state), and the solution of the stationary problem
of trans sonic flow around a solid is obtained as a limit
solution of the non-stationary problem as the time
tends to infinity. This method is very laborious, but is
quite feasible if electronic computers are employed.

As regards mathematical methods of study, super-
sonic aerodynamics may be subdivided into three fields.
1) Pure supersonic flow.
2) Mixed flow, with formation of local subsonic
zones.
3) High-temperature hypersonic flow producing
chemical reactions in the gas.
At higher velocities and, consequently, at higher tem-
peratures ionization takes place and radiation processes
become significant. This part of hypersonic aerodynam-
ics may be said to form a separate field, since the
mathematical problems involved are quite different
from those arising in the case of 'transparent' gases.
Problems in pure supersonic flow have been most
thoroughly investigated. This included the develop-
ment of numerical methods of characteristics, finite
difference methods and methods of semi-characteristics,
by which relatively simple solutions may be obtained
not only for plane-parallel or axi-symmetric but also
for spatial flows. The linear theory of supersonic flow
has also been exhaustively studied. Analytic solutions
can be obtained of numerous practical problems. The
latter may be somewhat complicated by the occurrence
of weak shock waves, but even these difficulties are of a
computational and not of a principle character.
Practically speaking, these problems are limited to
flows around sharply pointed solids and to the prob-
lems of internal aerodynamics (calculations of jet noz-
zles). However, if the supersonic velocities are very
high, sharp-edged solids are not employed (charring of
the edges), while a local subsonic flow always occurs at
the blunt end. Having computed this region, the
remaining (pure supersonic) flows are calculated by the
methods of supersonic aerodynamics.
As compared with the aerodynamics of transsonic
flow, the advantage of the problems of aerodynamics of
supersonic flow with local subsonic zones (even though
such zones are mixed) consists in the fact that such
zones usually represent bounded narrow domains in the
vicinity of the blunt ends. This is why it was possible to
develop effective numerical methods for the computa-
tion of local subsonic zones (the method of integral
relations, the method of inverse problems and the
adjustment method). It should be noted, however, that
up to the time of writing (1977) no rigorous
mathematical investigations of these problems have
ever been performed. There is also no proof that their
solutions do exists and are unique. Accordingly, numer-
ical methods are developed on the assumption that the
conditions of continuous velocities and accelerations
during the transition from the subsonic to the super-
sonic zone ensure the uniqueness of a solutions, which
is physically realistic. This hypothesis is confirmed by
all the results of numerical calculations.
Up till here, it has been said that the computation of
AERODYNAMICS, MATHEMATICAL PROBLEMS OF
the local subsonic zone (including the supersonic part
of the region of influence) is followed by calculations of
the flow by methods of purely supersonic aerodynam-
lCS.
There is another type of problems, which are con-
nected with hypersonic flow around thin solids (prior
to the initiation of chemical reactions). Hypersonic
flows around thin solids (excepting the blunt-shaped
regions) are distinguished by small changes in the velo-
city component along the main stream. For this reason
the equations may be simplified so that the problem of
flow around a thin body of a given shape (in the planar
and axi-symmetric cases) becomes identical to the one-
dimensional stationary problem. Many important qual-
itative characteristics of hypersonic flows have been
obtained in this way and approximate similitude laws
have been established; these are also extensively
employed in the analysis of the results of numerical cal-
culations, which are then reduced to very compact rela-
tions in a very wide range of Mach numbers and
geometric parameters of the solid.
Aerodynamic problems concerned with gases under-
going chemical reactions involve the simultaneous solu-
tion of the equations of motion and the equations of
chemical kinetics. While the resulting systems of equa-
tions are more complicated, the numerical methods
employed in their solution are essentially the same as
those employed in the aerodynamics of an ideal gas,
except that they are much more laborious. Such calcu-
lations are extensively carried out in practical work,
and since modelling natural conditions in aerodynamic
tubes cannot be realized in this temperature range, the
aerodynamic characteristics of hypersonic instruments
are mainly obtained by numerical methods.
There are two principal trends in the theory of
viscous fluids - the theory of the complete equations
of a viscous fluid (Navier - Stokes equations), and
boundary-layer theory. The boundary-layer equations
represent the principal term of the asymptotic expan-
sion of the Navier-Stokes equations in the vicinity of
the boundary of the solid to which the fluid particles
adhere, completely or in part. It is this fact which radi-
cally disturbs the solutions obtained for ideal fluids
near the surface of the solid around which the flow
takes place. The order of the error involved in the
boundary-layer equations is 1 / 5e, which means
that the theory is valid for large Reynolds numbers
only, and even then only in the regions of smooth con-
tinuous flow. Despite the fact that the boundary-layer
equations contain all the main terms of viscous stresses,
their mathematical structure is much simpler. While the
complete equations are of elliptic type, the boundary-
layer equations are parabolic with characteristics
directed along the normal to the surface of the solid. It
is therefore possible to conduct the computations 'layer
53
AERODYNAMICS, MATHEMATICAL PROBLEMS OF
by layer', i.e. to pass from one section of the boundary
layer to another, irrespective of the conditions prevail-
ing outside the range between the sections.
One consequence of the practical importance of the
theory of boundary layers (computations of the resis-
tance, of the surface temperatures, of the rate of des-
troyment of the surface under hypersonic flight condi-
tions) was the development of numerous approximate
methods for performing such calculations (the
Pohlhausen method, the one-parameter
Kochin - LoitsyanskiI method, etc.). However, if elec-
tronic computers are employed, all these approximate
methods become superfluous, since an accurate numeri-
cal solution of the boundary-layer equations presents
no difficulties, even in the difficult case of a hot gas in
which chemical reactions are taking place. The numeri-
cal routines may be so constructed that at each step the
system of equations is split, with respect to the x-
variable (along the tangent to the surface of the solid),
into a number of separate second-order differential
equations which, from the computational point of view,
is highly advantageous.
What was said above applies to plane-parallel and to
axi-symmetric flows. The structure of three-
dimensional boundary layers is more complicated. The
equations of the three-dimensional layers themselves
vary significantly with the geometry of the solid around
which the flow takes place. The development of
methods of calculation of three-dimensional boundary
layers is much less advanced, partly because calcula-
tions on plane sections are often used in practice; this,
despite the lack of a theoretical basis, often yields suffi-
ciently accurate results. There are no principal difficul-
ties in solving the exact equations of a three-
dimensional layer.
Except for a few special cases, in which analytic solu-
tions could be obtained, the development of methods
for solving the complete system of equations for a
viscous liquid is due to the introduction of electronic
computers. The principal difficulty in solving such
equations is their high order and the fact that the
region of influence is unbounded. In the simplest case
of stationary motion of a viscous incompressible fluid
with a constant viscosity coefficient one obtains a sys-
tem of two equations of the second order with respect
to the stream function t/; and vortex velocity w. They
may be written in the dimension-less form as follows:
~w = 2 [2f~-~~J
ay ax ax ay ,
~Ij; = w.
The boundary conditions on solid boundaries are only
for the stream function t/; = const, at/; / an = O. The solu-
tion of this non-linear system of the fourth order can
54
only be obtained by some kind of iterative process. The
iterative process is so constructed that separate equa-
tions for the vortex velocity and for the flow function
are obtained at each iteration step. Such a separation
must take place not only in the equations, but in the
boundary conditions as well. Thus, by choosing the
solution of the ideal fluid equations as the first
approximation, one obtains the following set of itera-
tions:
A _ aWn +1 _ [a<pn aWn _ a<pn aWn J
I.1Wn +1 2 ax - 2 ay ax ax ay ,
~o/n+l = Wn+h
with the following boundary conditions on the solid
boundaries:
alj;n
o/n+l = const, Wn+l = aa;-+wn'
where CP=t/;-y, which means that a separate set of
equations for the vortex flow and the flow function is
thus obtained in the (n + 1)-st approximation. The solu-
tion may be obtained in a similar manner by the
adjustment method, but it is not necessary in such a
case to take a real system of non-stationary equations;
rather, one may consider the system of parabolic equa-
tions: J
~
at = ~W-2[~~-~~
ay ax ax ay ,
~
at = ~,,-w
't',
which is more convenient to handle in numerical calcu-
lations.
The problem of flow around a body by a viscous
fluid involves not merely the development of computa-
tion methods, but also research on fundamental ques-
tions. Thus, it is unclear whether a stationary solution
exists at large Reynolds numbers (e.g. for the case of
flow around a cylinder). Numerical methods yield
acceptable solutions up to Reynolds numbers of several
hundreds, but above these values of the Reynolds
number the numerical calculations no longer converge.
Problems in the aerodynamics of radiant gases and
of rarefied gases have been solved only for the simplest
cases. However, in view of the practical importance of
these problems (in both these fields), mathematical
models of the relevant phenomena have been proposed.
Such models very often are no limiting cases of the
complete equations (model of gray radiation, model of
diffusive radiation in the theory of radiant gases,
Crook's model in the theory of rarefied gases), but
merely qualitative reflections of the relationships
described in the complete equations. The volume of the
computations involved in solving these problems is very
large, and success in this field will depend on the possi-
bilities for increasing the effectiveness of electronic
computers.

References
[1] KocHIN, N.E., KmEL', LA. and ROZE, N.V.: Theoretical hydro-
dynamics, Interscience, 1964 (translated from the Russian).
[2] KHRISTIANOVICH, S.A.: Gas flow around bodies with large sub-
sonic speed, Moscow, 1940 (in Russian).
[3] SEDOV, L.J. [L.I. SEDOV]: Two-dimensional problems in hydro-
dynamics and aerodynamics, Acad. Press, 1965 (translated from
the Russian).
[4] HAYES, W.: Hypersonic flow theory, Acad. Press, 1966.
[5] CHERNY!, 0.0.: Gas flow with large supersonic speed, Moscow,
1959 (in Russian).
[6] BELOTSERKOvsKIi, O.M. and CHusHKIN, P.I.: 'The numerical
solution of problems in gas dynamics', in Basic developments in
fluid dynamics, Vol. 1, Acad. Press, 1965, pp. 1-126.
[7] CHusHKIN, P.I.: The method of characteristics for spatial super-
sonic flows, Moscow, 1968 (in Russian).
[8] Flow around blunted bodies of a supersoniC gas flow, Moscow,
1967 (in Russian).
[9] Numerical methods in fluid mechanics, Wroclaw, 1969 (in Pol-
ish).
A.A. Dorodnitsyn
Editorial comments. In the equation of the conservation
of energy (g, V) denotes the scalar product of two vectors.
For some material on the branch or aerodynamics dealing
with such things as plasmas and the mechanics of stellar
gases (magneto- and electro-aerodynamics) see the article
Magneto-hydrodynamics, mathematical problems in.
References
[A1] BACHELOR, O.K.: An introduction to fluid dynamics, Cam-
bridge Univ. Press, 1970.
[A2] PAl, S.I.: Viscous flow, v. Nostrand, 1956.
[A3] HOWARTH, L. (ED.): Modern development in fluid dynamics.
High speed flow, 1-2, Oxford Univ. Press, 1953.
[A4] BlRKHOFF, G.: Hydrodynamics, a study in logic, fact and
similitude, Princeton Univ. Press, 1960.
[A5] KARMAN, TH. VON, ET AL. (EDS.): High speed aerodynamics
and jet propulsion, 1-2, Princeton Univ. Press, 1964.
[A6] PAl, S.I.: Radiation gas dynamics, Springer, 1966.
[A7] FERRARO, V.C.A. and PLUMPTON, c.: An introduction to
magneto-fluid mechanics, Oxford Univ. Press, 1966.
[A8] KOGAN, M.V.: Rarified gas dynamics, Acad. Press, 1969
(translated from the Russian).
[A9] BERS, L.: Mathematical aspects of transsonic gas dynamics,
Wiley, 1958.
AMS 1980 Subject Classification: 73JXX, 76GXX
AFFINE ALGEBRAIC SET, affine algebraic k-set -
The set of solutions of a given sy~em of algebraic
equations. Let k be a field and let k be its algebraic
closure. A subset X of the Cartesian product 1(n is said
to be an affine algebraic k-set if its points are the com-
mon zeros of some family S of the ring of polynomials
k[T]=k[TI' ... ,Tn]. The set Wx of all polynomials
in k [TI' ... ,Tn] that vanish on X forms an ideal, the
so-called ideal of the affine algebraic k-set. The ideal Wx
coincides with the radical of the ideal I (S) generated
by the family S, i.e. with the set of polynomials
fEk[TJ, ... ,Tn] such that fmEI(S) for some natural
number m (Hilbert's Nullstellensatz; cf. Hilbert theorem
3)). Two affine algebraic sets X and Y coincide if and
only if Wx = Wy. The affine algebraic set X can be
defined by a system of generators of Wx. In particular,
AFFINE CONNECTION
any affine algebraic set can be defined by a finite
number of polynomials (fl, ... ,fi)Ek[T]. The equal-
ities fl = ... = fi =0 are called the equations of X. The
affine algebraic sets of 'C form a lattice with respect to
the operations of intersection and union. The ideal of
the intersection X n Y is identical with the sum of
their ideals Wx + Wy , while the ideal of the union
X U Y is identical with the intersection of their ideals
Wx n Wy. Any set 1(n is an affine algebraic set, called
an affine space over k and denoted by AZ; to it
corresponds the zero ideal. The 'empty subset of In is
also an affine algebraic set with the unit ideal. The
quotient ring k[X]=k[TJ/ Wx is called the coordinate
ring of X. It is identical with the rin~ of k-regular func-
tions <,:n X, i.e. with the ring of k-valued functions,
f: X~k, for which there exists a polynomial FEk[T]
such thatf(x)=F(x) for all XEX. An affine algebraic
set is said to be irreducible if it is not the union of two
affine algebraic proper subsets. An equivalent defini-
tion is that the ideal Wx is prime. Irreducible affine
algebraic sets together with projective algebraic sets
were the subjects of classical algebraic geometry. They
were called, respectively, affine algebraic varieties and
projective algebraic varieties over the field k (or k-
varieties). Affine algebraic sets have the structure of a
topological space. The affine algebraic subsets are the
closed sets of this topology (the Zariski topology). An
affine algebraic set is irreducible if and only if it is
irreducible as a topological space. Further development
of the concept of an affine algebraic set leads to the
concepts of an affine variety and an affine scheme.
References
[1] ZARISKI, O. and SAMUEL, P.: Commutative algebra, 2, Springer,
1975.
[2] SHAFAREVICH, I.R.: Basic algebraic geometry, Springer, 1977
(translated from the Russian).
[3] HARTSHORNE, R.: Algebraic geometry, Springer, 1977.
1. V. Dolgachev
V.A. Iskovskikh
Editorial comments. A topological space is irreducible if
it is not the union of two closed proper subspaces.
AMS 1980 Subject Classification: 14A25, 14A 10,
14A05
AFFINE CONNECTION - A differential-geometric
structure on a smooth manifold M, a special kind of
connection on a manifold (cf. Connections on a mani-
fold), when the smooth fibre bundle E attached to M
has the affine space An of dimension n = dim M as its
typical fibre. The structure of such an E involves the
assignment to each point x EM of a copy of the affine
space (An)x, which is identified with the tangent
centro-affine space TAM). In an affine connection
each smooth curve L EM with origin Xo and each one
55
AFFINE CONNECTION

of its points XI is thus provided with an affine mapping A smooth field of frames is chosen in a coordinate
(An)x, ~(An)xo which satisfies the condition formulated neighbourhood of the origin Xo of the curve L, and the
below. Let M be covered with coordinate domains, image of the frame at point XI is defined as the solution
each provided with a smooth field of affine frames in {x(t), ei(t)} of the system
(An)x. The origin of these frames coincides with X (i.e. n
smooth vector fields, linearly independent at each point du = (W')X(I}(x(tu" 1 (5)
X of the domain, are given). The requirement is that, as du) = (wj)X(t)(x(tu"
t~O, when XI moves along L towards Xo, the mapping
for the initial conditions u(O)=O, ui(O)=ei' where
(An)x, ~(An)xo tends to become the identity mapping,
Xi =Xi(t) are the defining equations of the curve L. The
and that the principal part of its deviation from the
curve which is described in (An)x o by the point with
identity mapping be defined, with respect to some
position vector X (t) with respect to Xo is known as the
frame, by the system of linear differential forms
development of L. The field of frames in the coordinate
w' = ndx k , det I n 1=1= O,} neighbourhood may be so chosen that Wi =dx i ; then
wj = fJkwk. (1)
w; = r)k dxk. In the intersection of the coordinate
Thus, for (An)x, ~(An)xo' the image of the frame at XI is neighbourhoods, dx i =(OXi /oxj)wI, i.e.
the system consisting of the point in (An)x o with posi- Aj =ox i /ox j and
tion vector ei[wi(X)t +t:i(t)] and n vectors r L
,
, = axi ax] axk,
--~~f' +
2 x' ax' a (6)
e;[8j+wj(X)t+t:j(t)], where Xis the tangent vector to L ]k ax' ax) axk ]k ax) axk ax"
at Xo, and
Sjk = fjk-n),
af'1 ar 'k
1
a;' + r~kr'i/- r~/r'ik'
(7)
lim M = 0, lim (j(t) = 0.
t~O t t~O t Rjkl = ax~ -
A manifold M with an affine connection defined on Here Sjk and Rjkl are, respectively, the torsion tensor
it is called a space with an affine connection. During the and the curvature tensor of the affine connection on M.
transformation of a frame of the field at an arbitrary An affine connection on M may be defined by a system
point, XEM according to the formulas ei =A{-ej, of functions fjk on each coordinate neighbourhood
ej = A j ei, i.e. when passing to an arbitrary element of which transforms in the intersection of two neighbour-
the principal fibre bundle P of frames in the tangent hoods according to formula (5). The system fjk is
spaces (An)x with origins at the point x, the forms (1) called the object of the affine connection. The mapping
are replaced by the following I-forms on P: (An)x, is obtained with the aid of (5) into which
w' = Ajw', } w' = dx', wj = rjk dXk,
wi,] = AidA k + A i A I, wk (2)
is to be substituted.
k] k] I,

while the 2-forms If, in some neighbourhood of the point Xo, a vector
field X=~iei is given, then, when (An)x,~(An)xo' the
Q', = dw ,+W ,I\w' , }
l ]I
vector Xx (I) is mapped into the vector ~i(xl)ei(t) (where
(3)
Qj = dwj+wkl\wj, {ei(t)) is the solution of system (5. The differential of
are transformed as follows: this in (An)x o at t =0:

where Qi and Qj are composed from the forms (2)

(d~'+~'wj)e, = [aa;~ +~rjk] dxke"
according to (3). The equations (3) are called the struc- is called the covariant differential of the field X with
ture equations of the affine connection on M. Here the respect to the given affine connection. Here
left-hand sides - the so-called torsion forms Qi and
curvature forms Qj - are semi-basic (cf. Torsion form; \l t'
kS
= ~+tif'
axk""]k
Curvature form), i.e. they are linear combinations of the
form a tensor field, called the covariant derivative of the
wk Aw':
QI = .lSI
2 ] k w'/\ W k , 1 field X=~iei' If a second vector field Y=1{ek is given,
the covariant derivative of the field X in the direction
(4)
i_I I k 1
of Y is defined as
Q) - '2R]kIW I\w.

All I-forms Wi and w;,

defined on P and satisfying
which may also be defined with respect to an arbitrary
equations (3) with left-hand sides of type (4), define a
certain affine connection on M. The mapping field of frames by the formula
(An)x, ~(An)xo for a curve L EM is obtained as follows. w'(\l yX) = Y W'(X)+ w,,(Y)d(X).

56
 AFFINE COORDINATE FRAME

An affine connection on M may also be defined as a The concept of an affine connection arose in 1917 in
bilinear operator V which assigns a vector field V y X Riemannian geometry (in the form of the Levi-Civita
to each two vector fields X and Y, and which possesses connection); it found an independent meaning in 1918 -
the properties: 1924 owing to work of H. Weyl [1] and E. Cartan [2].
\7 y(jX) = (Yj)X + f\7 yX, \7 jYx = f\7 yX, References
[1) WEYL, H.: Raum, Zeit, Materie, Springer, 1923.
where f is a smooth function on M. The relation [2A) CARTAN, E.: 'Sur les varietes a connexion affine et la throrie
between these definitions is established by the formula de la reIativite generaIisee (premiere partie)', 40 (1923), 325-
V ek ej = f)k ej where {ej} is the field of frames. The
412.
[2B) CARTAN, E.: 'Sur les varietes a connexion affine et la theorie
fields of the torsion tensor and curvature tensor de la relativite generaIisee (premiere partie suite)" 41 (1924),
1-25.
S(X, Y) = S)kI'!T/ei = Oi(X, Y)ei' [2C] CARTAN, E.: 'Sur les varietes a connexion affine et la throrie
R(X, Y)Z = (R)k/~rI')I!ei = O)(X, y)wi(Z)ei de la relativite generaIisee (deuxieme partie)', 42 (1925), 17-88.
[3A) CARTAN, E.: 'Sur les varietes it connexion projective', Bull. Soc.
are defined by the formulas: Math. France 52 (1924), 205-241.
[3B) CARTAN, E.: 'Sur les espaces it connexion conforme', Ann. Soc.
S(X, Y) = \7x Y -\7yX-[X, Y], Polon. Math. 2 (1923), 171-221.
R(X, Y) = \7x\7yZ-\7y\7xZ -\7[X,y]z. [4) RASHEWSKI, P.K. [P.K. RAsHEVSKI11: Riemannsche Geometrie
und Tensorana/yse, Deutsch. Verlag Wissenschaft., 1959
A vector field X is said to be parallel along the curve L (translated from the Russian).
if V x(t)Xx(t) =0 holds identically with respect to t, i.e. [5) POSTNIKOV, M.M.: The variational theory of geodeSiCS, Moscow,
1965 (in Russian).
if, along L, (j Lumiste
d~i +I'!wj = 0,
Editorial comments. Instead of the articles [3A], [3B],
Parallel vector fields are used to effect parallel displace- one may consult [A 1]. Useful additional up-to-date refer-
ment of vectors (and, generally, of tensors) in an affine ences in English are [A2] and [A3].
connection, representing a linear mapping of the References
tangent vector spaces Tx,(M)~Txo(M), defined by the [A 1] LICHNEROWICZ, A.: Theorie globale des connexions et des
groupes d'holonomie, Cremonese, 1955.
mapping (An)x, ~(An)xo' In this sense any affine con-
[A2] KOBAYASHI, S. and NOMIzu, K.: Foundations of differential
nection generates a linear connection on M. geometry. 1, Wiley (Interscience), 1963.
A curve L is called a geodesic line in a given affine [A3] STERNBERG, S.: Lectures on differential geometry, Prentice-
Hall, 1964.
connection if its development is a straight line; in other
words, if, by a suitable parametrization, its tangent vec- AMS 1980 Subject Classification: 53B05
tor field x(t) is parallel to it. Geodesic lines are AFFINE COORDINATE FRAME - A set of n
defined with respect to a local coordinate system by the linearly-independent vectors ej (i = 1, ... ,n) of n-
system dimensionru affine space An, and a point O. The point
d 2 Xi dxi dx k
o is called the initial point, while the vectors ej are the
.
-2-+ f j kd - - d = O.
dt t t
scale vectors. Any point M is defined with respect to
Through each point, in each direction passes one geo- the affine coordinate frame by n numbers - coordi-
desic line. nates x j , occuring in the decomposition of the position
There is a one-to-one correspondence between affine vector OM by the scale vectors: OM=xses (summation
connections on M and connections in principal fibre convention). The specification of two affine coordinate
bundles of free affine frames in (An)x, x EM, generated frames defines a unique affine transformation of the
by them. To closed curves with origin and end at x space An which converts the first frame into the second
there correspond affine transformations (An)x ~(An)x' (see also Affine coordinate system).
which form the non-homogenous holonomy group of A.P. Shirokov
the given affine connection. The corresponding linear Editorial comments. An equivalent, and more usual,
definition is as follows. An affine coordinate frame in affine
automorphisms TX<M)~TX<M) form the homogeneous
n-space is a set of n + 1 points Po, ... ,Pn which are
holonomy group. In accordance with the holonomy
linearily independent in the affine sense, i.e. the vectors
theorem, the Lie algebras of these groups are defined by
Po Pi , i = 1, ... ,n, are linearly independent in the
the 2-forms of torsion OJ and curvature OJ. The Bianchi corresponding vector space. Independence of the vectors e;
identities apply to the latter: in the definition should be understood as independence in a
dO i = Oil\w
]
i -(.lI\W
J'
dO;J = Oik I\wJk -Wik I\ok
.OJ'
corresponding vector space.

In particular, for torsion-free affine connections, when References

[A1) SNAPPER, E. and TROYER, R.J.: Metric affine geometry,
OJ =0, these identities reduce to the following: Academic Press, 1971.

AMS 1980 Subject Classification: 51 N10

57
AFFINE COORDINATE SYSTEM
AFFINE COORDINATE SYSTEM - A rectilinear
coordinate system in an atTme space. An affine coordi-
nate system on a plane is defined by an ordered pair of
non-collinear vectors e) and e2 (an affine basis) and a
point 0 (the coordinate origin). The straight lines pass-
ing through the point 0 and parallel to the basis vec-
tors are known as the coordinate axes. The vectors e)
and e2 define the positive direction on the coordinate
axes. The axis parallel to the vector e) is called the
abscissa (axis), while that parallel to the vector e2 is
called the ordinate (axis). The affine coordinates of a
point M are given by an ordered pair of numbers (x, y)
which are the coefficients of the decomposition of the
vector OM by the basis vectors:
OM = xel +ye2.
The first number x is called the abscissa, while the
second number y is called the ordinate of M.
An affine coordinate system in three-dimensional
space is defined as an ordered triplet of linearly-
independent vectors e), e2, e3 and a point O. As in the
case of the plane, one defines the coordinate axes -
abscissa, ordinate and applicate - and the coordinates
of a point. Planes passing through pairs of coordinate
axes are known as coordinate planes.
A.S. Parkhomenko
AMS 1980 Subject Classification: 51 N1 0
AFFINE CURVATURE - A differential invariant of a
plane curve in the geometry of the general aff'me group
or a subgroup of it. The affine curvature is usually
understood to mean the differential invariant of the
curve in the geometry of the unimodular affine (or
equi-affine) group. In this geometry the affine (or, more
exactly, the equi-affine) curvature of a plane curve
y = y (x) is calculated by the formula
k = -t[(Y")-2/3f,
while the affine (or, more exactly, equi-affine) arc
length of the curve is
s = f(y'')1/3 dx.
There is a geometrical interpretation of the affine cur-
vature at a point M 0 of the curve: Let M be a point on
the curve close to M 0, let s be the affine length of the
arc MoM and let (J be the affine length of the arc of
the parabola tangent to this curve at M 0 and M. The
affine curvature at M 0 then is
ko = lim "' /nO(a-s).
M~Mo VS5
In the affine theory of space curves and surfaces there
are also notions of affine curvature which resemble the
respective notions of Euclidean differential geometry.
58
For references, see Affine differential geometry.
A.P. Shirokov
AMS 1980 Subject Classification: 53A15
AFFINE DIFFERENTIAL GEOMETRY - The branch
of geometry dealing with the differential-geometric pro-
perties of curves and surfaces that are invariant under
transformations of the affine group or its subgroups.
The differential geometry of equi-affine space has been
most thoroughly studied.
In an equi-affine plane any two vectors a, b have an
invariant (a, b) - the surface area of the parallelogram
constructed on a and b. With the aid of this concept,
the invariant parameter
s = jl (r,r) 11/3dt,
10
known as the equi-affine arc length, can be constructed
for a non-rectilinear curve r=r(t). The differential
invariant
k = [dds 2r ' .!!!.!.]
ds
2
3
is called the equi-affine curvature of the plane curve.
Constant equi-affine curvature characterizes curves of
the second order. A natural equation k = f (s) deter-
mines a curve up to an equi-affine transformation. The
vector n =d 2 r / ds 2 is directed along the affine normal
to a plane curve; the affine normal at a point M, k=j=O,
is the tangent to the locus of the mid-chords of the
curve parallel to the tangent at M, and coincides with
the diameter of the parabola which has third-order con-
tact with the curve at M.
Passing to the general affine group, two more invari-
ants of the curve are considered: the affine arc length a
and the affine curvature K. They can be expressed in
terms of the invariants s and k introduced above:
a = jkl/2ds K = _l_.dk
, k 3 / 2 ds'
(In equi-affine geometry, the magnitudes sand k them-
selves are called the affine arc length and the affine
curvature, for the sake of brevity.) The centro-affine arc
length, centro-affine curvature, equi-centro-affine arc
length and equi-centro-affine curvature of a plane curve
are constructed in a similar manner.
In equi-affine space it is possible to assign to any
three vectors a, b, c the invariant (a, b, c), which is the
volume of the oriented parallelepiped defined by these
vectors. The natural parameter (equi-affine arc length)
of a curve r=r(t) (rEe 3) is defined by the formula
s = j I (r, r, r) 11/ 6 dt.
10
The differential invariants K=(r', r"', r"''),
T= -(r", r''', r"''), where the primes denote differentia-
 AFFINE GEOMETRY

tion with respect to the natural parameter, are called,
respectively, the equi-affine curvature and the equi-affine
torsion of the spatial curve. The study of the curve is
reduced to selecting some moving frame; the frame
formed by the vectors
Two such forms, which satisfy supplementary differen-
tial conditions, determine the surface up to equi-affine
transformations. All these statements have appropriate
generalizations in the multi-dimensional case.
Many specific classes of surfaces are distinguished in

{r, r, 'r+K! } affine and equi-affine spaces: affine spheres (for which
the affine normals form a bundle), affine surfaces of
revolution (the affine normals intersect one proper or
and defined by the fourth-order differential neighbour-
improper straight line), affine minimal surfaces, etc.
hood of the curve being studied, is especially impor-
In addition to curves and surfaces, other geometrical
tant. The centro-affine theory of spatial curves has been
objects in equi-affine space are. also studied, such as
developed [5].
congruences and complexes of straight lines, vector
The following tensor is constructed for a non-
fields, etc.
developable surface r=r(u 1 , u 2) in equi-affine space:
In parallel with equi-affine differential geometry,
_ aij
development is also in progress of the differential
gij - laI 1 / 4 '
geometry of the general affine group and of its other
where aij =(rJ, r2, rij), a =det(ai), ri = air, rij = aijr. The subgroups both in three-dimensional and in multi-
vector dimensional spaces (centro-affine, equi-centro-affine,
affine-symplectic, bi-affine, etc.).
where V k is the symbol of the covariant derivative with References
respect to the metric tensor gij' determines the direction [I] BLASCHKE, W.: Vorlesungen aber Differentialgeometrie und
geometrische Grundlagen von Einsteins Relativitatstheorie. Affine
of the affine normal to the surface. The affine normal Differentialgeometrie, 2, Springer, 1923.
passes through the centre of the osculating Lie quadric. [2] SALKOWSKI, E.: Affine Differentialgeometrie, W. de Gruyter,
The derivational equations 1934.
IS [3] NORDEN, A.P.: Spaces with an affine connection, Nauka,
ajri = fijr, +gijN Moscow-Leningrad, 1976 (in Russian).
1k [4] LAPTEV, G.F.: 'Differential geometry of multi-dimensional sur-
define an intrinsic connection of the first kind fij of faces', Itogi Nauk. Geom 1963 (1965), 3-64 (in Russian).
[5] SHIROKOV, P.A. and SHIROKOV, A.P.: Differentialgeometrie,
the surface. There also arises at the same time an
2k Teubner, 1962 (translated from the Russian). A P Sh k
.. Iro ov
intrinsic connection of the second kind f ij , defined by
the derivational equations Editorial comments. For the development of affine dif-
2' ferential geometry after W. Blaschke, see [A1].
a/Ii = fijJl, + AijJl, References
where JI is a covariant vector defining the tangent plane [A1] SIMON, U.: 'Zur Entwicklung der affine Differentialgeometrie
nach Blaschke', in Wilhelm Blaschke gesammelte Werke, Vol.
to the surface and subject to the normalization condi-
4, Thales Verlag, 1985, pp. 35-88.
tion NJI= 1. The connections
lk 2k AMS 1980 Subject Classification: 53A 15
fij and fij
AFFINE DISTANCE - An invariant determined by
are conjugate with respect to the tensor gij in the sense
two line elements in an equi-aff'me plane. A point M
of A.P. Norden [3]. The tensor
together with a straight line m passing through it is
TkIJ = .12 [V-g]
IJ
called a line element (M, m). For two line elements
(M, m) and (N, n) the affine distance is 2j1 /3, where f
Ij'

which also plays a major part in projective differential is the surface area of the triangle MNP and P is the
geometry, makes it possible to construct the symmetric
point of intersection of the straight lines m and n. The
covariant tensor
Tijk = gksT'fj . affine distance for two elements tangent to a parabola
The two principal surface forms are also constructed: is equal to the affine arc length of this parabola (cf.
the quadratic form Affine parameter). In the three-dimensional equi-affine
space the affine distance may also be defined in terms
q, = gij du i du j of elements consisting of pairwise incident points,
and the Fubini - Pick cubic form straight lines and planes.
A.P. Shirokov
'" = Tijk du i du j du k .
AMS 1980 Subject Classification: 51 N1 0, 51 N25
These forms are connected by the apolarity condition
gijT;jk = o. AFFINE GEOMETRY - The branch of geometry

59
AFFINE GEOMETRY
whose subject are the properties of figures that are
invariant with respect to affine transformations (cf.
Affine transfonnation). Examples are the simple rela-
tion for three points to lie on a straight line, or the
parallelity of straight lines (planes). A.F. Mobius, in the
first half of the 19-th century, was the first to study the
properties of geometric images that go over into each
other as a result of affine transformations; however, the
concept of 'affine geometry' only arose following the
appearance in 1872 of the Erlangen program, according
to which each group of transformations has its own
geometry, the subject of which are the properties of the
figures that are invariant with respect to the transfor-
mations of this group. The group of affine transforma-
tions contains various subgroups; as a result, in addi-
tion to general affine geometry, subordinated
geometries equi-affine geometry, centro-affine
geometry, etc. - arose, corresponding to these groups.
Affine geometry also deals with problems of differential
geometry corresponding to specific transformation sub-
groups (cf. Affine differential geometry).
References
[1] ALEKSANDROV, P.S.: Lectures on analytical geometry, Moscow,
1968 (in Russian).
[2] EFiMOV, N.V.: Hohere Geometrie, Deutsch. Verlag Wissen-
schaft., 1960 (translated from the Russian).
E. V. Shikin
Editorial comments. Two English-language references
are [A 1], [A2].
References
[A 1] BORSUK, K.: Multidimensional analytic geometry, PWN, 1969.
[A2] MESERVE, B.E.: Fundamental concepts of geometry,
Addison-Wesley, 1955.
AMS 1980 Subject Classification: 51 N 10
AFFINE GROUP - The fundamental group of
transformations of an affine space. It is a subgroup of
the projective group and is represented by those projec-
tive transformations that map a fixed hyperplane of the
projective space into itself (cf. Projective transfonna-
tion).
A.P. Shirokov
AMS 1980 Subject Classification: 51 N1 0
AFFINE HULL of a set M in a vector space - The
intersection of all translated subspaces containing M.
V.A. Zalgaller
AMS 1980 Subject Classification: 15A03, 51 N1 0
AFFINE MINIMAL SURFACE - A surface whose
affine curvature is zero. As distinct from ordinary
minimal surfaces, which consist of saddle points only,
an affine minimal surface may contain elliptic points as
well. Thus, an elliptic paraboloid consists only of ellip-
tic points, and is an affine minimal surface.
E. V. Shikin
AMS 1980 Subject Classification: 53A 15
60
AFFINE MORPIDSM - A morphism of schemes
f : X ~S such that the pre-image of any open affine
subscheme in S is an affine scheme. The scheme X is
called an affine S-scheme.
Let S be a scheme, let A be a quasi-coherent sheaf of
(') s-algebras and let Ui be open affine subschemes in S
which form a covering of S. Then the glueing of the
affine schemes Spec f( Ui , A) determines an affine S-
scheme, denoted by SpecA. Conversely, any affine S-
scheme definable by an affine morphism f: X ~S is
isomorphic (as a scheme over S) to the scheme
SpecJ.( (') x). The set of S-morphisms of an S-scheme
f: Z ~S into the affine S-scheme SpecA is in bijective
correspondence with the homomorphisms of the
sheaves of (') s-algebras A ~J. ( (') z).
Closed imbeddings of schemes or arbitrary mor-
phisms of affine schemes are affine morphisms; other
examples of affine morphisms are entire morphisms
and finite morphisms. Thus the morphism of normali-
zation of a scheme is an affine morphism. Under com-
position and base change the property of a morphism
to be an affine morphism is preserved.
References
[1] GROTHENDIECK, A.: 'The cohomology theory of abstract alge-
braic varieties', in Proc. Internat. Math. Congress Edinburgh,
1958, Cambridge Univ. Press, 1960, pp. 103-118.
[2] DIEUDONNE, J. and GROTHENDIECK, A.: 'Elements de
geometrie algebrique', Pool. Math. IHES 4 (I960).
v.I. Danilov
I. V. Dolgachev
Editorial comments. f: X~S is a finite morphism if there
exist a covering (Sa) of S by affine open subschemes such
that f- 1 (Sa) is affine for all IX and such that the ring Ba of
f- 1 (Sa) is finitely generated as a module over the ring Aa of
Sa The morphism is entire if Ba is entire over A a, i.e. if
every xEBa integral over A a , which means that it is a root
of a monic polynomial with coefficients in A a , or,
equivalently, if for each xEBa the module Aa[x] is a
finitely-generated module over Aa.
References
[A1] HARTSHORNE, R.: Algebraic geometry, Springer, 1977.
AMS 1980 Subject Classification: 14A 15, 14E40
AFFINE NORMAL - A straight line defined in an
affine-invariant manner at each point of a hypersurface
in an affine space with the aid of the third-order dif-
ferential neighbourhood of this hypersurface; it is
essential that the principal quadratic form of the hyper-
surface does not degenerate. The affine normal at a
point M of a plane curve coincides with the diameter of
the parabola that has third-order contact with the curve
at M. If tangent hyper-quadrics are employed, a similar
interpretation can be given to the affine normal to a
hypersurface. In particular, the affine normal of a
hyper-quadric coincides with its diameter.
A. P. Shirokov
AMS 1980 Subject Classification: 53A 15

AFFINE PARAMETER, affine arc length - A parame-
ter on a curve which is preserved under transformations
of the affine group, for the determination of which the
derivatives of the position vector of the curve of the
lowest order must be known. The best known is the
parameter which is invariant with respect to the equi-
affine transformations, i.e. with respect to the affine
unimodular group. For a plane curve r=r(t) this affine
parameter is computed by the formula:
s = 11 (r,r) 11/3dt,
to
where (r, r) is the skew product of the vectors rand r.
In particular, for the affine length of an arc MoM 1 of a
parabola, the affine p'arameter is s = 2l /3, where f is
the area of the triangle formed by the chord MoM 1
and the tangents to the parabola at the points M 0 and
MI' It is possible to introduce in a similar manner the
affine parameter of a space curve in the geometry of
the general affine group or anyone of its subgroups.
A,P' Shirokov
Editorial comments. The arc length given by the formula
above is sometimes refered to as the special affine arc
length.
The notion of an 'affine parameter' is also used in the
theory of geodesics. An affine parameter of a geodesic (an
autoparaliel curve) of an affine connection is a parametriza-
tion x(t) of the geodesic such that for the corresponding
covariant derivative \l the equation
\l ;(t)x(t) =0
is valid (cf. [A3]).
References
[A1] SPIVAK, M.: Differential geometry, 2, Publish or Perish, 1970.
[A2] BUCHIN, S.: Affine differential geometry, Science Press; Gor-
don and Breach, 1983.
[A3] O'NEILL, B.: Semi-Riemannian geometry, Academic Press,
1983.
AMS 1980 Subject Classification: 53A 15, 53B05
AFFINE PSEUDO-DISTANCE The number
p=(MM', t), equal to the modulus of the vector pro-
duct of the vectors MM' and t, where M' is an arbi-
trary point in an equi-affme plane, M is a point on a
plane curve r = r(s), s is the affine parameter of the
curve and t=dr / ds is the tangent vector at the point
M. This number p is called the affine pseudo-distance
from M' to M. If M' is held fixed, while M is moved
along the curve, the affine pseudo-distance from M' to
M will assume a stationary value if and only if M' lies
on the affine normal of the curve at M. An affine
pseudo-distance in an equi-affine space can be defined
in a similar manner for a given hypersurface.
A.P. Shirokov
AMS 1980 Subject Classification: 53A15, 51 N25
AFFINE SCHEME
AFFINE SCHEME - A generalization of the concept
of an affine variety, which plays the role of a local
object in the theory of schemes. Let A be a commuta-
tive ring with a unit. An affine scheme consists _of a
topological space SpecA and a sheaf of rings A on
SpecA. Here, SpecA is the set of all prime ideals of A
(called the points of the affine scheme), provided with
the Zariski topology (or, which is the same, with the
spectral topology), in which the basis of open sets is
constituted by the subsets D(f)={pESpecA: fflp},
where f runs t~ough the elements of A. The sheaf of
local ~gs A is defined by the condition
f(D(f),A) = A f , where Af is the localization of the
ring A with respect to the multiplicative system
{i}n;;;.O, cf. Localization in a commutative algebra.
Affine schemes were first introduced by A. Grothen-
dieck [1], who created the theory of schemes. A scheme
is a ringed space which is locally isomorphic to an
affine scheme.
An affine scheme SpecA is called Noetherian
(integral, reduced, normal, or regular, respectively) if the
ring A is Noetherian (integral, without nilpotents,
integrally closed, or regular, respectively). An affine
scheme is called connected (irreducible, discrete, or
quasi-compact, respectively) if the topological space
SpecA also has these properties. The space SpecA of
an affine scheme is always compact (and usually not
Hausdorff).
The affine schemes form a category if the morphisms
of these schemes, considered as locally ringed spaces,
are considered as morphisms of affine schemes. Each
homomorphism of rings </>; A ~B define_s a morphism of
affine schemes: (SpecB, B)~(SpecA, A), consisting of
the continuous mapping a</>: SpecB~SpecA
(a<f>(p)=</I- 1(p) for _p~SpecB)! and a homomorphism of
sheaves of rings </I: A ~a</l:(B), which transforms the
section a / f of the sheaf A over the set D (f) into the
section <f>(a) / <P(f). The morphisms of an arbitrary
scheme (X, f) x) into an affine scheme (SpecA, A)
(which are also called X-valued points of SpecA) are in
a one-to-one correspondence with the homomorphisms
of rings A~f(X, f) x); thus, the correspondence
A ~(SpecA, A) is a contravariant functor from the
category of commutative rings with a unit into the
category of affine schemes, which establishes an anti-
equivalence of these categories. In particular, in the
category of affine schemes there are finite direct sums
and fibre products, dual to the constructions of the
direct sum and the tensor product of rings. The mor-
phisms of affine schemes which correspond to surjective
homomorphisms of rings are called closed imbeddings of
affine schemes.
The most important examples of affine schemes are
61
AFFINE SCHEME
affine varieties; other examples are affine group
schemes (cf. Group scheme).
In a manner similar to the construction of the sheaf
A it is p~ssible to ~nstruct, for any A-module M, a
sheaf of A-modules M on SpecA for which
f(D(j), M) = Mf = M AAf .
Such sheaves are called quasi-coherent. The category of
A-modules is equival~ent to the category of quasi-
coherent sheaves of A-modules on SpecA; projective
modules correspond to locally free sheaves. The coho-
mology spaces of quasi-coherent sheaves over an affine
scheme are described by Serre's theorem:
Hq(SpecA, M) = 0 if q>O.
The converse of this theorem (Serre's criterion for affin-
ity) states that if (X, (!) x) is a compact separable
scheme, and if H1(X, F)=O for any quasi-coherent
sheaf of (!) x-modules F, then X is an affine scheme.
Other criteria for affinity also exist [1], [4].
References
[I] GROTHENDIECK, A. and DIEUDONNE, J.: 'Elements de
geometrie algebrique', Publ. Math. IHES 4 (1960).
[2] DIEUDONNE, J.: 'Algebraic geometry', Adv. in Math. 1 (1969),
233-321.
[3] MANIN, Yu.I.: Lectures on algebraic geometry, 1, Moscow, 1970
(in Russian).
[4] GooDMAN, J. and HARTSHORNE, R.: 'Schemes with finite-
dimensional cohomology groups', Amer. J. Math. 91 (1969),
258-266.
v.I. Danilov
I. V. Dolgachev
Editorial comments. Reference [A1] is, of course, stan-
dard. It replaces [3]. An alternative to [1] is [A2].
References
[A1] HARTSHORNE, R.: Algebraic geometry, Springer, 1977,
[A2] GROTHENDIECK, A. and DIEUDONNE, J.: Elements de
geometrie algebrique, 1, Springer, 1971,
AMS 1980 Subject Classification: 14A15, 14L17,
14K99
AFFINE SPACE over a field k - A set A (the ele-
ments of which are called the points of the affine
space) to which corresponds a vector space Lover k
(which is called the space associated to A) and a map-
ping of the set A XA into the space L (the image of an
element (a, b)EA XA is denoted by ab and is called the
vector with beginning in a and end in b), which has the
following properties:
a) for any fixed point a the mapping x~a1, x EA, is
a bijection of A on L;
b) for any points a, b, c EA the relationship
ab+bt:+w = 0,
where -0 denotes the zero vector, is valid. The dimen-
sion of L is taken for the dimension of the affine space
A. A point a EA and a vector IE L define another
62
point, which is denoted by a + I, i.e. the additive group
of vectors of the space L acts freely and transitively on
the affine space corresponding to L.
Examples. 1) The set of the vectors of the space L is
the affine space A (L); the space associated to it con-
cides with L. In particular, the field of scalars is an
affine space of dimension 1. If L =k n, then A (kn) is
called the n-dimensional affine space over the field k, and
its points a=(a1,'" ,an) and b=(b 1, ... ,bn) deter-
mine the vector ab=(b 1 -a1,' .. ,bn-an)
2) The complement of any hyperplane in a projective
space over the field k is an affine space.
3) The set of solutions of a system of linear (alge-
braic or differential) equations is an affine space the
associated space of which is the space of solutions of
the corresponding homogeneous set of equations.
A subset A' of an affine space A is called an affine
subspace (or a linear manifold) in A if the set of vectors
ab, a, b EA " forms a subspace of L. Each affine sub-
space A' CA has the form a+L' ={a+/: IEL'}, where
L' is some subspace in L, while a is an arbitrary ele-
ment of A'.
A mapping J :A 1 ~A 2 between affine spaces A 1 and
A 2 is called affine if there exists a linear mapping of
the associated vector spaces cp: L 1~L2 such that
J(a+I)=J(a)+<p(I) for all aEA 1, IEL 1. A bijective
affine mapping is called an affine isomorphism. All
affine spaces of the same dimension are mutually iso-
morphic.
The affine isomorphisms of an affine space A into
itself form a group, called the affine group of the affine
space A and denoted by Aff(A). The affine group of
the affine space A (kn) is denoted by Affn(k). Each ele-
ment JEAffn(k) is given by a formula
fu], ... ,an = (b l , ,bn),
where
b, = ~a1aj +c;,
j
(a{) being an invertible matrix. The affine group Aff(A)
contains an invariant subgroup, called the subgroup oj
(parallel) translations, consisting of the mappings
J: A ~A for which cp: L~L is the identity. This group
is isomorphic to the additive group of the vector space
L. The mapping J~CP defines a surjective homomor-
phism of Aff(A) into the general linear group GL, with
the subgroup of parallel translations as kernel. If L is a
Euclidean space, the pre-image of the orthogonal group
is called the subgroup oj Euclidean motions. The pre-
image of the special linear group SGL is called the
equi-affine subgroup (d. Affine unimodular group). The
subgroup Ga CAffn(A) consisting of the mappings
f:A~A such that J(a +l)=a+<p(l) for a given aEA
and arbitrary I EL is called the centro-affine subgroup; it

is isomorphic to the general linear group GL of the
space L.
In algebraic geometry an affine algebraic set is some-
times called an affine space. A finite-dimensional affine
space can be provided with the structure of an affine
variety with the Zariski toplogy (cf. also Affine scheme).
Affine spaces associated with a vector space over a
skew-field k are constructed in a similar manner.
References
[I] BOURBAKI, N.: Elements 0fmathematics. Algebra: Algebraic
structures. Linear algebra, I, Addison-Wesley, 1974,
Chapt.l; 2 (translated from the French).
l. V. Dolgachev
A.P. Shirokov
Editorial comments. An affine isomorphism is also
called an affine collineation. An equi-affine group is also
called a Euclidean group.
References
[A1] BERGER, M.: Geometry, 1, Springer, 1987, Chap!. 2
(translated from the French).
AMS 1980 Subject Classification: 51 N1 0, 15A03
AFFINE SPHERE - A surface the affine normals (cf.
Affine normal) of which intersect at one point. An
affine sphere is said to be improper if this point is infin-
itely distant (or improper (ideal; otherwise it is proper.
In particular, any surface of the second order is an
affine sphere. An improper convex affine sphere is an
elliptic paraboloid.
E. V. Shikin
AMS 1980 Subject Classification: 53A 15
AFFINE TENSOR - An element of the tensor pro-
duct of p copies of an n-dimensional vector space En
and q copies of the dual vector space E~. Such a tensor
is said to be of type (p, q), the number p + q defining
the valency, or degree, of the tensor. Having chosen a
basis {e;} in En, one defines an affine tensor of type
(p, q) with the aid of nP +q components 1)', '..... J, which
transform as a result of a change of basis e; = A Jes
according to the formula
r", ...... ~ =A'I, "'A'ipAt, "'A/qT.' 'Sp
} 1}q SI Sp JI Jq tI ... 1'1 '
where AjA; =8). It is usually said that the tensor com-
ponents undergo a contravariant transformation with
respect to the upper indices, and a covariant transfor-
mation with respect to the lower.
A.P. Shirokov
Editorial comments. An affine tensor as described
above is commonly called Simply a tensor.
References
[A 1] DUBROVIN, B.A., FOMENKO, A.T. and NOVIKOV, S.T.: Modern
geometry - methods and applications, Springer, 1984
(translated from the Russian).
[A2] GREUB, W.H.: Multilinear algebra, Springer, 1967.
[A3] DODSON, c.T.J. and POSTON, T.: Tensor geometry, Pitman,
1977.
AMS 1980 Subject Classification: 15A69, 53A45
AFFINE TRANSFORMA nON
AFFINE TORSION - One of the differential invari-
ants of a curve in equi-affine space:
T = - (r", r"', r"''),
where the primes denote differentiation with respect to
the affine parameter of the position vector of the curve.
Similar concepts are introduced for curves in spaces
with other fundamental groups such as centro-affine.
A.P. Shirokov
Editorial comments.
References
[1] BUCHIN, S.: Affine differential geometry, Science Press; Gor-
don and Breach, 1983.
AMS 1980 Subject Classification: 53A 15
AFFINE TRANSFORMATION of a Euclidean space -
A one-to-one point mapping of the plane or space into
itself, where to three points lying on a straight line
correspond three points also lying on a straight line.
Thus, by an affine transformation straight lines are
transformed into straight lines. An affine transforma-
tion of the plane transforms intersecting lines into
intersecting lines, and parallel lines into parallel lines.
An affine transformation of space results in an affine
mapping of each plane onto some plane; intersecting
planes are mapped into intersecting planes, while paral-
lel planes are mapped into parallel planes. Moreover,
the mutual location of two straight lines is preserved:
intersecting straight lines are mapped into intersecting
straight lines, parallel straight lines are mapped into
parallel straight lines, while skew lines are mapped into
skew lines.
Under an affine transformation the ratio between
directed segments lying on the same line or on parallel
straight lines is equal to the ratio of their images. The
ratio between the surface areas of two measurable fig-
ures (on the Euclidean plane) and the ratio between the
volumes of two measurable bodies (in the Euclidean
space) are preserved as well. Under an affine transfor-
mation a set of vectors in the plane (in space) is one-
to-one mapped on a set of vectors in the plane (in
space), and this mapping is linear. An affine transfor-
mation is defined in an affine coordinate system by a
non-degenerate (non-homogeneous) linear transforma-
tion; thus, in the case of a plane, an affine transforma-
tion is analytically expressed by the formulas
x = CIIX +CIV +C13,
y' = C21 X + C2V+ C23,
with the supplementary condition
Affine transformations in space are defined in a similar
manner.
63
AFFINE TRANSFORMATION
Under an affine transformation an algebraic curve is
converted into an algebraic curve, the order of the
curves being preserved. In particular, a curve of the
second order is converted into a curve of the second
order, moreover, ellipses are converted into ellipses,
hyperbolas into hyperbolas, parabolas into parabolas,
etc.
Examples of affine transformations: isometric
transformations, similarity transformations, and uni-
form contractions of a plane to a straight line. Each
affine transformation in a plane is a product of an
isometric transformation and two uniform contractions
to two mutually-perpendicular straight lines. Each
affine transformation in a space is a product of an
isometric transformation and three uniform contrac-
tions to three pairwise-perpendicular planes.
The affine transformations form a group; the similar-
ity transformations constitute a subgroup of this group;
the set of isometric transformations is a subgroup of
the group of similarity transformations.
Affine transformations are the most general one-to-
one mappings of a plane (a space) into itself under
which straight lines are preserved.
References
(1) ALEKSANDROV, P.S.: Lectures on analytical geometry, Moscow,
1968 (in Russian).
[2) POSTNIKOV, M.M.: AnalytiC geometry, Moscow, 1973 (in Rus-
sian).
A.S. Parkhomenko
Editorial comments. In terms of coordinates like above
similarity transformations take the form
x' = C11X+C12Y+C13, Y' = -ec12x+ec11y+c23
Similitudes, or homothetic transformations, which in terms of
coordinates take the form x' = C11 x+ C13, i
= C22Y+ C23'
constitute another interesting subgroup of the group of
affine transformations.
References
[A 1) COXETER, H.S.M.: Introduction to geometry, Wiley, 1961.
[A2] MESERVE, B.E.: Fundamental concepts of geometry,
Addison-Wesley, 1955.
AMS 1980 Subject Classification: 51 N1 0
AFFINE UNIMODULAR GROUP, equi-affine group -
The subgroup of the general affine group consisting of
the affine transformations of the n-dimensional affine
space
X 1-+ X = Ax+a (*)
that satisfy the condition det(A) = 1. If the vectors x
and x are interpreted as rectangular coordinates of
points in the n-dimensional Euclidean space En, then
the transformation (*) will preserve the volumes of n-
dimensional domains of P. This makes it possible to
introduce the concept of volume in an equi-affine
space, which is a space with a fundamental affine uni-
modular group. If, in formulas (*), one puts a = 0, then
64
one obtains a centro-affine unimodular group of
transformations isomorphic to the group of all matrices
of order n with determinant equal to one. Such a
group of matrices is called the unimodular group or spe-
ciallinear group of order n and is denoted by SL(n).
A.P. Shirokov
AMS 1980 Subject Classification: 20H05
AFFINE VARIETY, affine algebraic variety - A gen-
eralization of the concept of an affme algebraic set. An
affine variety is a reduced affme scheme X of finite
type over a field k, i.e. X = SpecA, where A is a com-
mutative k-algebra of finite type without nilpotent ele-
ments. The affine variety X=Speck[T 1 , . . , Tn],
where k [T 1, . . . , Tn] is the ring of polynomials over k,
is called affine space over k and is denoted by AZ. An
affine scheme is an affine variety if and only if it is iso-
morphic to a reduced closed subscheme of an affine
space. Each system of generators Xl, . . . ,Xn of a k-
algebra A defines a surjective homomorphism
cp: k[TJ, ... ,Tn]~A, defined by the formula CP(Ti)=Xi'
Let "k be the algebraic closure of k. The subset of the
set "k n consisting of the common zeros of all the poly-
nomials of the ideal ker cp is an affine algebraic set over
k. The coordinate ring of such an affine algebraic set is
isomorphic to the ring A. Each affine algebraic set over
k in turn defines an algebraic variety Speck[X], where
k[X] is the coordinate ring of X. The set of points of
an affine variety is in a one-to-one correspondence with
the irreducible subvarieties of the corresponding affine
algebraic set.
To each affine variety X=SpecA a functor on the
category of k-algebras is assigned. It is defined by the
correspondence:
B ~ X(B) = Homk -alg(A, B).
If ~=k (respectively, if B=k), the elements of the set
X(k) (respectively, X(k)) are called ge<!..metric (respec-
tively, rational) points of X. The set X (k) is in a bijec-
tive correspondence with the set of maximal ideals
Specm (A) of the ring A, and with the set of points of
an algebraic set V whose coordinate ring is isomorphic
to A. The spectral topology in the space X induces on
the everywhere-dense subset Specm (A ) a topology
which corresponds to the Zariski topology on V.
I. V. Dolgachev
Editorial comments. Frequently the name 'variety'
means a reduced and irreducible scheme of finite type over
an algebraically closed field.
References
[A 1] SHAFAREVICH, I.R.: BaSIC algebraic geometry, Springer, 1977
(translated from the Russian).
AMS 1980 Subject Classification: 14A 10
 AIRY FUNCTIONS

AFFINITY - A brief synonym for an affine transfor-
mation, defined as a non-degenerate linear substitution
X' = A~xs +a'.
A.P. Shirokov
AMS 1980 Subject Classification: 51 N10
AFFINOR - An affine tensor of type (1, 1). Specify-
ing an affinor with components f} is equivalent to
specifying an endomorphism of the vector space
according to the rule vi =/1 vS To the iden ti ty
endomorphism there corresponds a unique affinor. The
correspondence by which the matrix If} I is assigned
to each affinor realizes an isomorphism between the
algebra of affinors and the algebra of matrices. An affi-
nor is sometimes defined in the literature as a general
(affine) tensor.
A.P. Shirokov
Editorial comments. I.e. one is concerned here with the
isomorphism VI8i V~End(V) of linear algebra.
so are w( wz) and w( w2 z), where w = e 2m j3, and any
two of these solutions are linearly independent. The fol-
lowing identity holds:
w(z)+w(wz)+w(w 2z) _ O.
References
[1) AIRY, G.B.: 'On the intensity of light in the neighbourhood of
a caustic', Trans. Cambridge Phi/os. Soc. 6 (1838), 379-402.
[2) BABICH, V.M. and BULDYREv, V.S.: Asymptotic methods in
problems of diffraction of short waves, Moscow, 1972 (in Rus-
sian).
[3) ABRAMOWITZ, M. and STEGUN, l.A. (EDS.): Handbook of
mathematical functions, Appl. Math. Series, Vol. 55, Nat.
Bureau of Standards, 1964.
M. V. Fedoryuk
AMS 1980 Subject Classification: 34A30
AIRY FUNCflONS - Particular solutions of the Airy
equation.
The first Airy function (or simply the Airy function)
is defined by

Icos [3-T +
AMS 1980 Subject Classification: 15A69, 53A45
Ai(x) = -;1
00
]
xt dt.
AFFIx OF A COMPLEX NUMBER Z = a + bi in its
geometric representation - The point in the complex For complex values of z
plane corresponding to this number, i.e. the point with
Cartesian coordinates (a, b). The affix is sometimes 1 [exp [
Ai(z) = 2", zt-t3 ] dt,
3

identified with the complex number itself.

where y = (ooe -2?Ti /3,0] U [0, + 00) is a contour in the
AMS 1980 Subject Classification: 30-01
complex t-plane. The second Airy function is defined
AIRY EQUATION - The second-order linear ordi- by
Bi(z) = iw
2 Ai(w2Z)-iwAi(wz), w 27Ti 3 =e / .
nary differential equation
The functions Ai(x) and Bi(x) are real for real x.
y" -xy = O. A second collection of Airy functions was introduced
It occurred first in G.B. Airy's research in optics [1]. Its by V.A. Fock [V.A. Fok]:
general solution can be expressed in terms of Bessel
functions of order +- I / 3: v(z) = V; Ai(z),

w,(z) = 2e i7T / 6v(wz),

y(x) = C,Vxl'/3 [fiX3/2]+C2Vxl_,/3 [fiX3/2].
W2(Z) = 2e- i7T / 6v(w-'z);
Since the Airy equation plays an important role in vari- in this case v(z) is called the Airy- Fok function
ous problems of physics, mechanics and asymptotic (Airy- Fock function). The following identities hold:
analysis, its solutions are regarded as forming a distinct
w,(Z)-W2(Z) --
class of special functions (see Airy functions). v(z) = 2i ,w,(z) = W2(Z)' (1)
The solutions of the Airy equation in the complex
plane z, Any two of v(z), wJ(z) and W2(Z) are linearly indepen-
w"-zw = 0, dent.
have the following fundamental properties: The most important Airy function is v(z) (or Ai(z.
I) Every solution is an entire function of z and can Its asymptotic behaviour on the real axis is given by
be expanded in a power series
vex) = ~ X-.;!;4 exp[-fx3/2][l+O(x-3/2)], x~+oo,
w(z) = w(o) (l+~+ Z6
23 (2,3)'(5'6)
+ ... ] +
v(x) = IX~/4 [Sin[tlxI3/2+~]+O(IXI-3/2)J.
+w'(O) [z+ ;.: + (3.4;.:6'7) + ... ], x~-oo,

which converges for all z.
2) If w(z) 'jE is a solution of the Airy equation, then
so v(x) decreases rapidly for x >0, x I and oscillates
strongly for x<O, I x I 1. The functions Wl(X) and

65
AIRY FUNCTIONS

W2(X) increase exponentially as x~+oo. For complex where q(x) is a smooth real-valued function on the
z the Airy functions have the following asymptotic interval I=[a, b] and A>O is a large parameter. The
expansions as I z I~ 00: zeros of q(x) are called turning points (or transfer

v(z) '" 2~z-'/4exp [_f 3/2] x

z (2)
points) of the equation (3). Let
a < Xo < b, q(xo) = 0, q' (xo) =1=
(such a point is called simpfe),
for I argz I,,;;;;,w-,
00
x ~ (-Iranz -3n /2
n =0
q(x) =1= for xEI, x=l=xo, q'(xo) > O.
Set

00
X ~ anz -3n /2
n =0
for I argz - -
W

3
I ";;;;'W-, ((x) ~ [~!. Vq(t) dtr" ';80 ((x) ~ ';8o(x - x,>
where Yo(x) = ((x-'/2Ai(-A2 / 3x,
_ r [3n+t]9- n Y,(x) = ((x-'/2Bi(-A2 / 3X.

an - (2n)! Equation (3) has linearly independent solutions Yo(x)

andYI(x) such that, as A~+oo,
The asymptotic expansion of W2(Z) is of the form (2),
but it is valid in the sector
Yi x ) = Yi x ) [1+0 [~ J]. a";;;;'x";;;;'xo, j=O, I,

iarg [z+;] I,,;;;;, W-.

yo(x) = Yo(x) [I+O[ ~]]+Y,(x)o[f],

Here E(O,77) is arbitrary, the branches of and Vz
z 1 / 4 are positive on the semi-axis (0, 00), and the
asymptotic expansions are uniform with respect to argz y,(x) = Y,(x) [I +0 [~]] + Yo(x)O [fJ,
and can be differentiated term by term any number of
Xo ,,;;;;, x , ;;;;, b,
times. In the remaining sector I arg - z I < the asymp-
totic expansion of v(z) is expressed in terms of those uniformly with respect to x.
of WI (z) and W 2 (z) by means of (1); hence, the asymp- This result has been generalized in various directions:
totic expansion of v(z) has a different form in different asymptotic series have been obtained for the solutions,
sectors of the complex z-plane. This fact was first esta- the case q = q(x, A) has been studied (for example, if
blished by G.G. Stokes [2] and is called the Stokes q(x, A) can be expanded in an asymptotic series
phenomenon. q'" '2/::=OA -nqn(x) as A~+ 00), and the asymptotic
The Airy functions occur in the study of integrals of behaviour of the solutions near multiple turning points
rapidly-oscillating functions, of the form has been investigated. Other generalizations concern the
b
equation w" +A2 q(X)W = 0, (4)
leA, a) = jtfAS(x,a)f(x, a)dx,
a where the function q(z) is analytic in a domain D of
for A>O, A~oo. Here f and S are smooth functions, S the complex z-plane. Let f be the maximal connected
is real and 0: is a real parameter. If for small values of component of the level line
0:;;00 the phase S has two close non-degenerate station- Z

ary points XI(O:) and X2(0:) that coincide for 0:=0, for Re j v'q(t) dt = 0,
example, if Zo

emanating from a turning point z 0 and containing no

sex, a) = ax - x 3 + 0(X4) as x-)oO,
then for small values of 0:;;00, as A--+ + 00, the contribu-
tion to the asymptotics of the integral coming from a
neighbourhood of the point x = 0 can be expressed in
terms of the Airy function v and its derivative (see [6]).
Integrals of this kind occur in the study of short-wave
fields near a simple focus (see [7] and [8]); the Airy
functions arose in connection with the study of this
problem [1].
Consider the second-order differential equation
y" +A 2q(X)y = 0,
other turning points; then f is called a Stokes line. If
q = - z (that is, (4) is the Airy equation), then the
Stokes lines are the rays (- 00,0) and (0, ei1T /3).
Analogously, if zo is a simple turning point of (4), then
there are three Stokes lines fl' f2 and f3 emanating
from it and the angle between adjacent lines at z 0 is
equal to 277 /3. Let Si be a neighbourhood of z 0 from
which a neighbourhood of the Stokes line fi' j = 1, 2, 3,
has been removed. For a suitable numbering of the Si'
equation (4) has three solutions w/z), j = 1, 2, 3, such
(3) that, as ;\--+ + 00,

66

for Z ESj.
The Airy functions also occur in the study of asymp-
totic solutions of ordinary differential equations and
systems of higher order near simple turning points.
References
[I] AIRY, G.B.: Trans. Cambridge Phi/os. Soc. 6 (1838), 379-402.
[2] STOKES, G.G.: Trans. Cambridge Phi/os. Soc. 10 (1857), 105-
128.
[3] Fmc, V.A.: 'Tables of the Airy functions', Moscow, 1946
(in Russian).
[4] SOOUN, A. and ABRAMOWITZ, M.: Handbook of mathematical
functions, Appl. Math. Ser., Vol. 55, Nat. Bur. Stand., 1970.
[5] BABICH, V.M. and BULDYREv, V.S.: Asymptotic methods in
problems of diffraction of short waves, Moscow, 1972 (in Rus-
sian).
[6] FEDORYUK, M.V.: The saddle-point method, Moscow, 1977 (in
Russian).
[7] LANDAU, L.D. and LIFSHITS, E.M.: The classical theory of
fields, Addison-Wesley, 1951 (translated from the Russian).
[8] MASLOV, V.P. and FEDORYUK, M.V.: Quasi-classical approxi-
mation for the equations of quantum mechanics, Moscow, 1976
(in Russian).
[9] DORODNITSYN, A.A.: 'Asymptotic laws of distribution of the
characteristic values for certain types of second-order differen-
tial equations', Uspekhi Mat. Nauk 6, no. 7 (1952),3-96 (in
Russian).
[10] WASOW, W.: Asymptotic expansions for ordinary differential
equations, Interscience, 1965.
(II] FEDORYUK, M.V.: Asymptotic methods for linear ordinary dif-
ferential equations, Moscow, 1983 (in Russian).
M. V. Fedoryuk
Editorial comments. The Airy function can be expressed
in terms of modified Bessel functions of the third kind:
Ai(x) = ?T~ VxK, j3 [ ; X2/3].
The function Ai(z) satisfies the differential equation
w"(z)=zw(z), cf. [A2].
References
[A1] OLVER, F.W.J.: Asymptotics and special functions, Acad.
Press., 1974.
[A2] LEBEDEV, N.N.: Special functions and their applications,
Dover, reprint, 1972 (translated from the Russian).
AMS 1980 Subject Classification: 33A40
AITKEN SCHEME - A method for computing the
value at a point x of the interpolation polynomial
4(x) with respect to the nodes Xo, ... ,Xn, based on
the successive application of the formula
Lk(X) = L(o,. ,k)(X) = (*)
1 IL(O, ... ,k-l)(X) xo-XI
Xk-XO L (l, ... ,k) (x) xk-x,k=1,2, ... ,
where L(i, ... ,m)(x) is the interpolation polynomial with
interpolation nodes Xi, .. , X m , in particular,
L(i)(x) = f(Xi) (see Interpolation fonnula). The process
of computations by means of (*) may finish if the
ALBANESE VARIETY
values of two interpolation polynomials of consecutive
degrees coincide in the required number of decimal
places. The Aitken scheme is convenient for interpolat-
ing the values of a function given in the form of a table
(of values), by renumbering the interpolation nodes in
the order in which I x - Xi I increases.
References
[I] BEREZIN, I.S. and ZHIDKOV, N.P.: Computing methods, I, Per-
gamon Press, 1973 (translated from the Russian).
(2) BAKHVALov, N.S.: Numerical methods: analysis, algebra, ordi-
nary differential equations, Mir, 1977 (translated from the Rus-
sian).
M,K Samarin
Editorial comments. Aitken published the gist of his
method in [A1]. Aitken's scheme is disadvantageous when a
number of interpolations must be carried out over the same
range. In such cases, an alternative to the computation of
the Lagrange polynomials is the use of divided differences
in the construction of Newton's interpolation formula.
References
[A1] AITKEN, A.c.: 'On interpolation by iteration of proportional
parts, without the use of differences', Proc. Edingburgh
Math. Soc. 3, no. 2 (1932), 56-76.
[A2] HILDEBRAND, F.B.: Introduction to numerical analysis,
McGraw-Hili, 1974.
AMS 1980 Subject Classification: 65D05
ALBANESE VARIETY - An Abelian variety Alb(X)
canonically attached to an algebraic variety X, which is
the solution of the following universal problem: There
exists a morphism cJ>: X ~Alb(X) such that any mor-
phism f : X ~A into an A~elian variety A factors into a
product f= fcJ>, where f:A~Alb(X) (so named in
honour of G. Albanese). If X is a complete non-
singular variety over the field of complex numbers, the
Albanese variety can be described as follows. Let [21 be
the space of everywhere-regular differential forms of
degree I on X. Each one-dimensional cycle y of the
topological space X determines a linear function
1
W~ w on [21. The image of the mapping
y
HI(X, Z)~([2I) thus obtained is a lattice r in ([21)*,
and the quotient space ([21)* / r coincides with the
Albanese variety of X. From the algebraic point of
view, an Albanese variety may be considered as a
method of defining an algebraic structure on some quo-
tient group of the group Z of zero-dimensional cycles
of degree 0 on X. If X is a non-singular complete alge-
braic curve, both its Picard variety and its Albanese
variety are called its Jacobi variety. If the ground field
has characteristic zero, then the equalities
dimAlb(X) = dimk HO(X, Q.\-) = dimk H1(X, @ x)
are valid. The number dimAlb(X) is called the irregu-
larity of the variety X. If the field has finite characteris-
tic, the inequalities
irr(X) ,;;;;; dimHo(X, Q.\-) and irr(X)';;;;; dimH1(X, @ x)
67
ALBANESE VARIETY
hold. If the ground field has positive characteristics it
can happen that
dimHl(X, tJ x) *- dimHo(X, n.\}
The Albanese variety is dual to the Picard variety.
References
[1) BALDASSARRI, M.: Algebraic varieties, Springer, 1956.
[2) LANG, S.: Abelian varieties, Springer, 1983.
A.N. Parshin
AMS 1980 Subject Classification: 14K30, 14KXX
ALBEDO METHOD (in transport theory) - A
method of solving boundary value problems for a tran-
sport equation. The Albedo method is a variant of the
matrix factorization method, the role of factorization
coefficients being played by the matrices of reflection
and transmission of a sequence of layers of increasing
thickness. Practically speaking, the Albedo method is
used only in problems in one-dimensional spaces, both
to find the reflection and transmission coefficients and
to find the solution of the transport equation in a
medium.
References
[1) BEREZIN, I.S. and ZHIDKOY, N.P.: Computing methods, Per-
gamon Press, 1973 (translated from the Russian).
[2) GERGOMENOYA, T.A., ET AL.: Al'bedo neutrons, Moscow, 1973
(in Russian).
TA. Germogenova
AMS 1980 Subject Classification: 35-XX, 82A70
ALEKSANDROV-CECH HOMOLOGY AND
COHOMOLOGY, spectral homology and cohomology -
Homology and cohomology theories which satisfy all
Steenrod - Eilenberg axioms (except, possibly, the
exactness axiom) and a certain continuity condition.
The Aleksandrov-Cech homology groups (modules)
Hn(X, A ; G) [1], [2] are defined as the inverse limit
lim..... Hn(a, a'; G) over all open coverings a of the space
X; here a signifies not only the covering, but also its
nerve, and a' is the subcomplex in a that is the nerve
of the restriction of a on the closed set A (cf. Nerve of
a family of sets). The possibility of passing to the limit
is ensured by the existence of simplicial projections
(/3, /3')~(a, a') defined, up to an homotopy, by the
inclusion of {1 in a. The Aleksandrov - tech cohomology
groups Hn(x, A ; G) are defined as the direct limit
linL. Hn(a, a'; G). The homology groups satisfy all the
Steenrod - Eilenberg axioms except for the exactness
axiom. All axioms are valid for cohomology, and,
partly for this reason, cohomology is often more useful.
The exactness axiom is also valid for homology on the
category of compacta if G is a compact group or field.
In addition, Aleksandrov-Cech homology and coho-
mology groups have the property of continuity: For
X=lim.- X h the homology (cohomology) groups in
68
addition are equal to the respective limit of the homol-
ogy (cohomology) groups of the compacta X).. The
Aleksandrov-Cech theory is the only theory satisfying
the Steenrod - Eilenberg axioms (with the exception
indicated above) and this condition of continuity. On
the category of paracompact spaces, the usual charac-
terization by mappings into Eilenberg - MacLane
spaces is valid for cohomology; while the cohomology
itself is equivalent to the cohomology defined in sheaf
theory. Cohomology may also be defined as cohomol-
ogy of some cochain complex, which makes it possible
to operate with sheaves of cochains. Similar ideas,
applied to homology, are contained in the homology
theory originating from N. Steenrod, A. Borel and oth-
ers, which satisfies all axioms including the exactness
axiom (but ty.e property of continuity is lost).
Aleksandrov-Cech homology and cohomology, includ-
ing the above modification, are employed in homologi-
cal problems in the theory of continuous mappings, in
the theory of transformation groups (a connection with
quotient spaces), in the theory of generalized manifolds
(in particular, in various duality relations), in the
theory of analytic spaces (e.g. in defining the funda-
mental classes of homology), in homological dimension
theory, etc.
References
[1) ALEKSANDROFF, P.S. [P.S. Ai.EKSANDROY): 'Untersuchungen
tiber Gestalt und Lage abgeschlossener Mengen beliebiger
Dimension', Ann. of Math. (2) 30 (1929), 101-187.
[2) CECH, E.: 'Theorie genet-ale de l'homologie dans un espace
que1conque', Fund Math. 19 (1932), 149-183.
[3) STEENROD, N.E. and E1LENBERG, S.: Foundations of algebraiC
topology, Princeton Univ. Press, 1966.
[4) SKLYARENKO, E.G.: 'Homology theory and the exactness
axiom', Russian Math. Surveys 24, no. 5 (1969), 91-142.
(Uspekhi Mat. Nauk 24, no. 5 (1969), 87-140)
E. G. Sklyarenko
Editorial comments. Often one ~peaks also of Cech
cohomology instead of Aleksandrov-Cech cohomology.
AMS 1980 Subject Classification: 55N05
ALEKSANDROV COMPACflFICATION, Aleksan-
drov compact extension - The unique Hausdorff compac-
tifieation aX of a locally compact, non-compact, Haus-
dorff space X, obtained by adding a single point 00 to
X. An arbitrary neighbourhood of the point 00 must
then have the form {oo} U (X \ F), where F is some
compactum in X. The Aleksandrov compactification aX
is the smallest element in the set B (X) of all compactif-
ications of X. A smallest element in the set B (X) exists
only for a locally compact space X and must coincide
with ax'
The Aleksandrov compactification was defined by
P.S. Aleksandrov (1) and plays an important role in
topology. Thus, the Aleksandrov compactification aRn

of the n-dimensional Euclidean space is identical with
the n-dimensional sphere; the Aleksandrov compactifi-
cation aN of the set of natural numbers is
homeomorphic to the space of a convergent sequence
together with the limit point; the Aleksandrov compac-
tification of the 'open' Mobius strip coincides with the
real projective plane RP2. There are pathological situa-
tions connected with the Aleksandrov compactification,
e.g. there exists a perfectly-normal, locally compact and
countably-compact space X whose Aleksandrov com-
pactification has the dimensions dim aX < dim X and
Ind aX < Ind X.
References
[I] ALEKSANDROFF, P.S. [P.S. ALEKSANDROV]: 'Ueber die Metrisa-
tion der im Kleinen kompakten topologischen Raumen', Math.
Ann. 92 (1924), 294-301.
v. V. Fedorchuk
Editorial comments. The Aleksandrov compactification
is also called the one-point compactification.
References
[A1] DUGUNDJI, J.: Topology, Allyn and Bacon, 1966. Theorem
8.4.
AMS 1980 Subject Classification: 54035
ALEPH, ~ - The first letter of the Hebrew alphabet.
As symbols alephs were introduced by G. Cantor to
denote the cardinal numbers (the cardinality cf. Cardi-
nal number) of infinite well-ordered sets. Each cardinal
number is some aleph (a consequence of the axiom of
choice). However, many theorems about alephs are
demonstrated without recourse to the axiom of choice.
For each ordinal number a by ~a=w(wa) one denotes
the cardinality of the set of all ordinal numbers smaller
than Wa' In particular, ~o is the cardinality of the set of
all natural numbers, ~I is the cardinality of the set of
all countable ordinal numbers, etc. If a</3, then
~a<~fl' The cardinal number ~a+1 is the smallest
cardinal number which follows ~a' The generalized con-
tinuum hypothesis states that 21'1 = ~a + I for any ordinal
number a. If a = 0, the equation assumes the form
21'1 =~], and forms the content of the continuum
0 P.s. Aleksandrov
hypothesis. The set of all alephs smaller than ~a is
totally ordered according to magnitude, and its order
type is a. The definitions of the sum, the product and a
power of alephs are obvious. One has
Ka + Kil = Ka 'KIl = Kmax(a. Ill'
The following formulas are most-frequently encoun-
tered. The recursive Hausdorff formula:
a particular case of which, for a=O, is the Bernshtefn
formula:
The recursive formula of Tarski: If an ordinal number
ALEXANDER DUALITY
a is a limit ordinal, and if /3<cf(a), then
K!(p
a
= ~
~
K!(p
<.
a
Here cf(a) denotes the confinal character of the ordinal
number a. As in the case of cardinal numbers, one dis-
tinguishes between singular alephs, regular alephs, limit
alephs, weakly inaccessible alephs, strongly inaccessible
alephs, etc. For example, ~a is singular if a is a limit
ordinal and if cf(a)<a. There is no largest aleph
among all alephs. It was shown by Cantor that the set
of all alephs is meaningless, i.e. that there is no such
set. See also Totally weD-ordered set; Continuum
hypothesis; Set theory; Ordinal number; Cardinal
number.
References
[I] ALEKSANDROV, P.S.: EinfUhrung in die Mengenlehre und die
Theorie der reelen Funktionen, Deutsch. Verlag Wissenschaft.,
1956 (translated from the Russian).
[2] HAUSDORFF, F.: Set theory, Chelsea, reprint, 1978.
[3] COHEN, P.J.: Set theory and the continuum hypothesis, Benja-
min, 1966.
[4] KURATOWSKI, K. and MOSTOWSKI, A.: Set theory, North-
Holland, 1968.
B.A. Efimov
Editorial comments. A more recent theorem on the
exponentiation of alephs was proved by J. Silver in 1974, cf.
[A2]. A particular case says that if
2!( = K~+1 for all ~<W1
then
A reasonable up-tO-date additional reference for this topic is
[A 1].
References
[A1] LEVY, A.: Basic set theory, Springer, 1979.
[A2] SILVER, J.: 'On the singular cardinals problem', in Proc. Inter-
nat. Congress Mathematicians Vancouver, 1974, Vol. 1,
1975, pp. 265-268.
AMS 1980 Subject Classification: 03E55, 04A 10,
04A30
ALEPH-ZERO - A cardinal number which is the car-
dinality of all countably-infinite sets. Denoted by ~o.
AMS 1980 Subject Classification: 03E10, 04A 10
ALEXANDER DUALITY - The connection between
the homological properties of complementary subsets of
a topological space, owing to which the homological
properties of a set can be defined by certain properties
of its complement. The first theorems of this kind were
formulated in terms of set-theoretic rather than alge-
braic topology. It was shown by C. Jordan in 1892 that
a simple closed curve divides the plane into two
domains and forms the common boundary of the two
(the Jordan theorem). This theorem was extended in
1911 by H. Lebesgue and L.E.J. Brouwer (indepen-
69
ALEXANDER DUALITY
dently of each other) to the case of an n-dimensional
manifold in an (n + I)-dimensional spherical (or
Euclidean) space; a connection was also established
between this fact and the property of an r-dimensional
manifold (in an n-dimensional space) being linked with
an (n-r-l)-dimensional manifold (Lebesgue). It was
shown by Brouwer in 1913 that the number of domains
into which a plane is subdivided by a closed subset
depends only on the topological properties of this set.
This type of duality was first expressed by J.W. Alex-
ander in 1922 [1] in purely homological terms.
Alexander's theorem [2], [3], [4] states that the r-
dimensional Betti number mod 2 of a (finite)
polyhedron <P in an n-dimensional spherical space is
equal to the (n -r-l)-dimensional Betti number mod
2 of its complement. P.S. Aleksandrov in 1927 extended
this theorem to the case of any closed set <P. The dual-
ity expressed by this theorem is known as Alexander
duality.
The next important step in the development of dual-
ity of this kind was the theorem of Pontryagin [2], [3],
[4] (1934), which states that the r-dimensional homology
group H,(A, X) of a closed set A in an n-dimensional
spherical manifold M n with a compact group X of coef-
ficients, and the (n - r - 1)-dimensional homology
group Hn-,-J(B, Y) of the complement B=Mn\A
with the (discrete) group Y of coefficients which is dual
to X in the sense of the theory of characters, are duals,
and their scalar product defines the linking coefficient
of arbitrary cycles from the homology classes scalarly
multiplied. This theorem is known as the
Alexander- Pontryagin theorem; the duality formulated
in it is known as Alexander- Pontryagin duality or as
Pontryagin duality. A series of subsequent publications
led to Aleksandrov's theorem [5], [7], the formulation of
which differs from that of Pontryagin's theorem in that
A may be an arbitrary subset of M n , the group X may
be both compact or discrete, and H,(A, X) and
Hn -, -J (B, Y) are understood to mean the
Aleksandrov-Cech homology groups (cf.
Aleksandrov-Cech homology and cohomology), one of
which has compact support, while the other is of
spectral type. The forms of Alexander duality for arbi-
trary sets are obtained by replacing the latter type of
groups by their dual cohomology groups of the same
dimension over the dual group of coefficients.
In subsequent generalizations, spherical manifolds
were replaced by more general manifolds (homology
manifolds which are acyclic in certain dimensions, cf.
Homology manifold), the Aleksandrov-Cech groups
were replaced by Sitnikov - Steenrod groups (groups of
projective type) and other groups; the groups of coeffi-
cients were replaced by modules, sheaves, etc.
70
References
[1] ALExANDER, II, J.W.: 'A proof of the invariance of certain
constants of analysis situs', Trans. Amer. Math. Soc. 16 (1915),
148-154.
[2] ALEKsANDROV, P.S.: Combinatorial topology, Graylock, Roches-
ter, 1956 (translated from the Russian).
[3] PONTRYAGIN, L.S.: 'The general topological theorem of duality
for closed sets', Ann. of Math. 35, no. 4 (1934), 904-914.
[4] LEFSCHETZ, S.: Algebraic topology, Arner. Math. Soc., 1955.
[5] ALEKSANDROV, P.S.: 'Fundamental duality theorems for non-
closed sets of an n -dimensional space', Mat. Sb. 21 (63), no.
2 (1947), 161-232 (in Russian).
[6A] ALEKSANDROY, P.S.: 'Topological duality theorems I. Closed
sets', Trudy Mat. Inst. Steklov. 48 (1955) (in Russian).
[6B] ALEKSANDROY, P.S.: 'Topological duality theorems II. Non-
closed sets', Trudy Mat. Inst. Steklov. 54 (1959) (in Russian).
[7] CHOGOSHYILI, G.S.: 'Homology theory of compact spaces', in
Proc. 4-th. All-Union Math. Congress, Vol. 2, Moscow, 1964,
pp. 57-62 (in Russian).
[8] KAPLAN, S.: 'Homology properties of arbitrary subsets of
Euclidean spaces', Trans. Amer. Math. Soc. 62 (1947), 248-271.
[9] BOURGIN, D.G.: Modern algebraic topology, Macmillan, 1963.
[10] CHoGOSHYILI, G.S.: 'On homology theory of non-closed sets',
in General topology and its relations to modern analysis and
algebra. Proc. Symp. Prague, 1961, pp. 123-132.
G.S. Chogoshvili
AMS 1980 Subject Classification: 55NXX, 54FXX
ALEXANDER INVARIANTS - Invariants connected
with the module structure of the one-dimensional
homology of a manifold M, freely acted upon by a free
Abelian group Ja of rank a with a fixed system of gen-
erators t J, ,ta _
The projection of the manifold M onto the space M
of orbits (cf. Orbit) is a covering which corresponds to
the kernel Ka of the homomorphism y: G-,,;Ja of the
fundamental group 'TTJ(M)=G of the manifold M. Since
Ka='TTJCM), the group Ba=Ka/ Ka, where Ka is the
commutator subgroup of the kernel K a , is iSOlporphic to
the one-dimensional homology group H J (M, Z). The
extension l-,,;Ka-,,;G-,,;Ja-,,;I generates the extension
(*): I-,,;Ba-,,;G / Ka-,,;Ja-,,;I, which determines on Ba
the structure of a module over the integer group ring
Z(JG) of the group Ja (cf. Group algebra). The same
structure is induced on Ba by the given action of Ja on
M. Fixation of the generators ti in Ja identifies Z(JG)
with the ring
La=La(tJ, ... ,ta)=Z[tJ,t- J, ... ,ta,t;J] of Laurent
polynomials in the variables ti . Purely algebraically the
extension (*) defines and is defined by the extension of
modules (**): O-";Ba--,,Aa-,,;fa-";O [5]. Here fa is the ker-
nel of the homomorphism f.: La-";Z (f.tj = I). The
mod~le Aa is called the Alexander module of the cover-
ing M --"M. In the case first studied by J.W. Alexander
[1] when M=M(k) is the complementary space of
some link k of multiplicity JL in the three-dimensional
sphere S3, while the covering corresponds to the com-
mutation homomorphism Y/l: G(k)--"P of the link

group, A I' is the Alexander module of the link k. The
principal properties of G which are relevant to what
follows are: G / G' is a free Abelian group, the defect
of the group G is 1, G has the presentation
{XI, ... ,Xm+l;rJ, ... ,rm } for which Yp.(XJ=ti'
1~i ~JL; Y.(Xi) = 1, i>JL (cf. Knot and link diagrams). In
the case of links the generators ti EP correspond to the
meridians of the components k i C k and are fixed by
the orientations of these components and of the sphere.
As a rule, M is the complementary space M(k) of k,
consisting of JL (n - 2)-dimensional spheres k i in S n. In
addition to the homomorphism Ym' one also considers
the homomorphism yo: G(k)~J, where y(x) is equal to
the sum of the link coefficients of the loop representing
x with all ki .
The matrix [)(a of the module relations of a module
Aa is called the Alexander covering matrix and, in the
case of links, the Alexander link matrix. It may be
obtained as the matrix
[ ~]Y"<P,
aXj
where {Xi; rd is a presentation of the group G. If JL= 1,
the matrix ina of module relations for Ba is obtained
from [)(a by discarding the zero column. The matrices
[)(a and ina are defined by the modules Aa and Ba up
to transformations corresponding to transitions to other
presentations of the module. However, they can be used
to calculate a number of module invariants. Alexander
ideals are ideals of the module A a, i.e. series of ideals
Ei(Aa) of the ring
La: (O)C,E o C, ... c,Ei - 1 C,Eic' ... C,(l), where Ei is
generated by the minors of [)(a of order
(m-i)X(m-i) and Ei=La for m-i<1. The opposite
numbering sequence may also be employed. Since La is
both a Gaussian ring and a Noetherian ring, each ideal
Ei lies in a minimal principal ideal (Ll i); its generator fli
is defined up to unit divisors t~. The Laurent polyno-
mial Lli(t J, ... ,t1') is simply call~ the Alexander poly-
nomial of k (or of the covering M~M). If fli7'~::O, it is
multiplied by tt, , ... ,t~ so that fli(O, ... ,0)*0 and
*00. To the homomorphism YrJ there correspond a
module A, ideals Ei and polynomials ai, designated,
respectively, as Alexander's reduced module, Alexander's
reduced ideals and Alex..ander's reduced polynomials of.!
(or of the covering M rJ~M). If JL= 1, then A =A.
[)(A) is obtained from [)( by replacing all ti by t. If
JL~2, a l is divisible by (t -l)r2. The polynomial
V(t)=a l (t) / (t -1)1'-2 is known as the Hosokawa
polynomial. The module properties of A (k) have been
studied [4], [8], [10]. The case of links has not yet _been
thoroughly investigated. For JL=l, the group HI(M;R)
is finitely generated over any ring R containing Z in
ALEXANDER INVARIANTS
which Ll(O) is invertible [7], in particular over the field
of rational numbers, and, if fl(O) = + 1, then also over
Z. In such a case fl(t) is the cparacteristic_polynomial
of the transformation t:HI(M;R)~HI(M;R). The
degree of fll (t) is equal to the rank of HJ (M; R); in
particular, fll(t)=l if and only if HI(M;Z)=O. If
n = 3, the link ideals have the following symmetry pro-
perty: Ei =Ei' where the bar denotes that the image is
taken under the automorphism generated by replacing
all ti by t i- I. It follows that
A ( -I
I.li t I , ... , tI'-I) -- tN,
I ... tI'N I.li
A (
t J, . . ,t I' ) f or certaIn
.
integers N i This symmetry is the result of the
Fox - Trotter duality for knot and link groups. It may
also b~ deduced from the Poincare duality for the mani-
fold M, taking into account the free action of Ja [3]. If
fll (t I, ... ,t1')*0, then the chain complex C.(M) over
the field of fractions PI' of the ring L I' is acyclic
(n = 3), and the Reidemeister torsion TEP I' / II
corresponding to the imbedding Lp. cP 1" where II is
the group of units of Lp., is defined accordingly. If
JL=2, then T=fl l ; if JL= 1, then T=fll / t -1 (up to
units of Lp.). The symmetry of fll for n = 3 is a conse-
quence of the symmetry of T. If {t= 1, it follows from
the symmetry of fli (t) and from the property fli ( 1) = + 1
that the degree of fli(t) is even. The degree of V(t) is
also even [4]. The following properties of the knot poly-
nomials fli(t) are characteristic: Lli(l) = +1;
fli(t) = t 2k fli(t -I); fli+ I divides fl i ; and fl; = 1 for all i
greater than a certain value N, i.e. for each selection
fli(t) with these properties there exists a knot k for
which they serve as the Alexander polynomials. The
Hosokawa polynomials [4] are characterized by the
property V(t)=t 2k V(t-I) for any {t~2; the polynomi-
als Lli of two-dimensional knots by the property
fl l (1)=1.
Alexander invariants, and in the first place the poly-
nomials, are powerful tools for distinguishirJg knots and
links. Thus, Lli fails to distinguish between only three
pairs out of the knots in a table containing fewer than
9 double points (cf. Knot table). See also Knot theory;
Alternating knots and links.
References
[1) ALEXANDER, J.W.: 'Topological invariants of knots and links',
Trans. A mer. Math. Soc. 30 (1928), 275-306.
[2) REIDEMEISTER, K.: Knotentheorie, Chelsea, reprint, 1948.
[3) BLANCHFIELD, R.C.: 'Intersection theory of manifolds with
operators with applications to knot theory', Ann. of Math. (2)
65, no. 2 (1957), 340-356.
[4) HOSOKAWA, F.: 'On 'V -polynomials of links', Osaka J. Math.
10 (1958),273-282.
[5) CROWELL, R.H.: 'Corresponding groups and module
sequences', Nagoya Math. J. 19 (1961),27-40.
[6) CROWELL, R.H. and Fox, R.H.: Introduction to knot theory,
Ginn, 1963.
[7) NEUWIRTH, L.: Knot groups, Princeton Univ. Press, 1965.
71
ALEXANDER INVARIANTS
[8) CROWELL, R.H.: 'Torsion in link modules', J. Math. Mech. 14,
no. 2 (1965), 289-298.
[9) LEVINE, J.: 'A method for generating link polynomials', Amer.
J. Math. 89 (1967),69-84.
[10) MILNOR, J.W.: 'Multidimensional knots', in Conference on the
topology of manifolds, Vol. 13, Boston, 1968, pp. 115-133.
A. V. Chernavskif
AMS 1980 Subject Classification: 55NXX, 57TXX,
57MXX
AL'FA - An algorithmic language developed in the
USSR to solve scientific and technical problems with
the aid of computers. Al'fa is to some extent an exten-
sion of the reproducing properties of Algol. To the
quantities appearing in Algol, complex numbers and
component quantities with an internal dimension (vec-
tor, matrices, etc.) have been added. Other additions
include descriptions which serve to introduce designa-
tions for components of complex and multi-
dimensional quantities. A special class of function rou-
tines, the body of which is defined by an expression,
has been introduced. It is possible to display the count-
ing of natural numbers by dots, to use summation and
product signs and chains of inequalities. There is a sup-
plementary kind of loops in which the desired number
of repetitions can be performed without introducing an
enumeration parameter. The principal programming
systems include Al'fa (for M-20 type machines), Algibr
(for the complex M-220/BESM-6) and Al'fa-6 (for
BESM-6 machines).
References
[1) ERSHOV, A.P., KOZHUKHIN, G.L and VOLOSHIN, YU.M.:
Operating languages for systems of automatic programming,
Novosibirsk, 1964 (in Russian).
[2) Al'fa - a system for automatization ofprogramming, Novosi-
birsk, 1967 (in Russian).
(3) ERSHOV, A.P., KOZHUKHIN, G.L and POTIOSIN, LV.: Hand-
book for the use of the system Al'fa, Novosibirsk, 1968 (in Rus-
sian).
(4) Manual for the use of systems of automatization of programming,
Novosibirsk, 1975 (in Russian).
A.P. Ershov
AMS 1980 Subject Classification: 68B05
ALGAMS - An algorithmic language developed for
use mainly on medium-power electronic computers. It
was developed in 1963 - 1966 by the Group for
Automatic Programming for Medium-type Computers
(GAMS), appointed by the commission for comprehen-
sive collaboration of the acadeInies of sciences of
socialistic countries. It is supposed to be a standard
language for the exchange of algorithms between social-
istic countries. It is based on the language Algol-60 (d.
Algol), on which certain restrictions were imposed in
order to facilitate the translation process. The most
important restrictions include the prohibition of a
recursive use of procedures, the demand of a specific a-
72
tion of a procedure's formal parameters, the description
of identifiers (except for labels) prior to their use, and a
simplified structure of the naming expressions. These
restrictions are related to those imposed on Algol-60 in
the unified language known as the Algol-60 subset. At
the same time new features were introduced into
Algams, including both external and part identifiers,
which do not exist in Algol-60. External identifiers are
names of arrays recommended for inclusion in the
external memory of the machine. Reading and writing
of the external arrays is achieved by a standard
exchange procedure. Part identifiers, preceding the
arrays, isolate parts of the program which may be
stored in the external memory, and are recalled to the
working memory upon entry into the respective array.
In this way the effectiveness of the language is
enhanced, while the capacity of the working memory of
the machine remains relatively small. Moreover,
Algams contains detailed descriptions of input and out-
put procedures, as well as a more precise definition of
procedure bodies in the case of a description of them in
other languages.
References
(II 'Description of the language ALGAMS', in Algorithms and
algorithmic languages, Moscow, 1968, pp. 3-56 (in Russian).
[21 LYUBIMsKIi, E.Z. and MARTYNYUK, V.V.: 'Refinement of the
definition of the Algams language', Programming and Computer
Software 2 (1976), 73-74. (Programmirovanie 1 (1976), 87-88)
V. V. Lutsikovich
AMS 1980 Subject Classification: 68BXX, 6SF05
ALGEBRA - 1) A branch of mathematics (see Alge-
bra). The term is used in combinations such as homo-
logical algebra, commutative algebra, linear algebra,
multilinear algebra, and topological algebra.
2) A special case of an operator ring: An algebra
(occasionally, a linear or vector algebra) over a field,
over a skew-field or over a commutative ring. Associa-
tive algebras (formerly described as 'hypercomplex sys-
tems') and non-associative algebras are algebras in this
sense.
3) A synonym of universal algebra. This includes
algebras such as Boolean algebras, unary algebras, etc.
(d. Boolean algebra; Unary algebra).
AMS 1980 Subject Classification: OS-XX, 12-XX, 13-
XX, 15-XX, 16-XX, 17 -XX, 18-XX, 20-XX
ALGEBRA - The branch of mathematics in which
one studies algebraic operations (cf. Algebraic opera-
tion).
Historical survey. The simplest algebraic operations -
arithmetic operations on positive integers and positive
rational numbers - can be encountered in the oldest
mathematical texts, which indicates that the principal

properties of these operations were known even in early
antiquity. In particular, the Arithmetic of Diophantus
(third century A.D.) had a major influence on the
development of algebraic ideas and symbols. The term
'Algebra' originates from the work AI jabr al-muqabala
by Mohammed al-Khwarizmi (nineth century A.D.),
which describes general methods for solving problems
which can be reduced to algebraic equations of the first
and second degree. Towards the end of the fifteenth
century the cumbersome verbal descriptions of
mathematical operations which had previously pre-
vailed began to be replaced by the contemporary sym-
bols '+' and '-', and subsequently symbols for
powers, roots, and parentheses appeared. F. Viete, at
the end of the sixteenth century, was the first to use the
letters of the alphabet to denote the constants and the
variables in a problem. Most of the present-day sym-
bols of algebra were known as early as the mid-
seventeenth century, which marks the end of the
'prehistory' of algebra. The development of algebra
proper took place during the next three centuries, dur-
ing which views as to the proper subject matter of this
discipline kept radically changing.
In the seventeenth and eighteenth centuries 'Algebra'
was understood to mean the science of computations
carried out on algebraic symbols - 'identity' transfor-
mations of formulas consisting of letters, solving alge-
braic equations (cf. Algebraic equation), etc. - as dis-
tinct from arithmetic, which dealt with calculations per-
formed on explicit numbers. It was assumed, however,
that the symbols stood for actual numbers: integers or
fractions. A brief table of the contents of one of the
best textbooks of that time, L. Euler's Introduction to
algebra includes integers, ordinary and decimal frac-
tions, roots, logarithms, algebraic equations of degrees
one to four, progressions, additions, Newton's binomial
and Diophantine equations. Thus, by the mid-
eighteenth century, algebra corresponded, more or less,
to the 'elementary' algebra of our own days.
The principal subject dealt with by the algebra of the
eighteenth and nineteenth centuries were polynomials.
Historically, the first problem was the solution of alge-
braic equations in one unknown, i.e. equations of the
type:
The purpose was to derive formulas expressing the
roots of the equation in terms of its coefficients, by
means of addition, multiplication, subtraction, division
and extraction of roots ('solution by radicals').
Mathematicians were able to solve first- and second-
degree equations even in the earliest times. Substantial
advances were made in the sixteenth century by Italian
mathematicians: a formula was found for solving
third-degree equations (cf. Cardano fonnula) and
ALGEBRA
fourth-degree equations (cf. Ferrari method). During
the following three .::enturies fruitless efforts were made
to find similar formulas for solving equations of higher
degrees; in this connection, the problem of finding at
least a 'formula-free' proof of the existence of a com-
plex root of an arbitrary algebraic equation with com-
plex coefficients became of major interest. This theorem
was first stated in the seventeenth century by A.
Girard, but was rigorously proved by C.F. Gauss only
towards the end of the eighteenth century (cf. Algebra,
fundamental theorem of). Finally, it was established by
N.H. Abel in 1824 that equations of degree higher than
four cannot, in general, be solved by radicals, and E.
Galois in 1830 stated a general criterion for the solva-
bility of algebraic equations by radicals (cf. Galois
theory). Other problems were neglected at that time,
and algebra was understood to mean the 'analysis of
equations', as noted by J. Serret in his course of higher
algebra (1849).
Studies on algebraic equations in one unknown were
accompanied by studies on algebraic equations in
several unknowns, in particular of systems of linear
equations. The study of linear equations resulted in the
introduction of the concepts of a matrix and a deter-
minant. Matrices subsequently became the subject of an
independent theory, the algebra of matrices, and their
scope of application was extended beyond the solution
of systems of linear equations.
From the mid-nineteenth century onwards studies in
algebra gradually moved away from the theory of equa-
tions towards the study of arbitrary algebraic opera-
tions. The first attempts at an axiomatic study of alge-
braic operations dates back to the 'theory of relations'
of Euclid, but no progress was made in this direction,
since a geometrical interpretation of even the simplest
arithmetic operations - ratios of lengths or of areas -
is impossible. Further progress only became possible
as a result of gradual generalization and intensive study
of the concept of a number, and of the appearance of
arithmetic operations performed on objects entirely
unlike any number. The first such examples were
Gauss' 'composition of binary quadratic forms', and P.
Ruffini's and A.L. Cauchy's multiplication of permuta-
tions. The abstract concept of an algebraic operation
appeared in the mid-nineteenth century in the context
of studies on complex numbers (cf. Complex number).
There appeared G. Boole's algebra of logic, H.
Grassmann's exterior algebra, W. Hamilton's quatern-
ions (cf. Quaternion) and A. Cayley's matrix calculus,
while C. Jordan published a major treatise on pennuta-
tion groups.
These studies prepared the way for the transition of
algebra at the turn of the nineteenth century into its
modern stage of development, which is characterized by
73
ALGEBRA
the combination of previously separate algebraic ideas
on a common axiomatic basis and by a considerable
extension of the scope of its applications. The modern
view of algebra, and of the general theory of algebraic
operations, crystallized at the beginning of the twen-
tieth century under the influence of D. Hilbert, E.
Steinitz, E. Artin and E. Noether, and was fully esta-
blished by 1930 with the appearance of B.L. van der
Waerden's Modern algebra.
The subject matter of algebra, its principal branches and
its connection with other branches of mathematics. The
subject matter of modern algebra are sets and algebraic
operations on these sets (i.e. algebras or universal alge-
bras, cf. the terminology in Algebra; Universal algebra),
considered up to an isomorphism. This means that,
from the point of view of algebra, the sets themselves
and the sets as carriers of algebraic operations are
indistinguishable, and in this sense the proper subject
of study are the algebraic operations themselves.
For a long time the studies that were actually carried
out concerned only a few basic types of universal alge-
bras which naturally appeared in the development of
mathematics and its applications.
One of the most important and most thoroughly stu-
died type of algebras is a group, i.e. an algebra with
one associative binary operation, containing a unit ele-
ment and, for each element, an inverse element. The
concept of a group was, historically, the first example
of a universal algebra and in fact served, in many
respects, as a model for the construction of algebra and
of mathematics in general at the turn of the nineteenth
century. Independent studies on generalizations of
groups such as semi-groups, quasi-groups and loops
began only much later (d. Loop; Quasi-group; Semi-
group).
Rings and fields are very important types of algebras
with two binary operations. The operations in rings and
fields are usually called addition and multiplication. A
ring is defined by Abelian group axioms for the addi-
tion, and by distributive laws for multiplication with
respect to addition (d. Rings and algebras). Originally,
only rings with associativity multiplication were stu-
died, and this requirement of associativeness occasion-
ally forms part of the definition of a ring (d. Associa-
tive rings and algebras). The study of non-associative
rings (cf. Non-associative rings and algebras) is today a
fully recognized independent discipline. A skew-field is
an associative ring in which the set of all non-zero ele-
ments is a multiplicative group. A field is a skew-field
in which multiplication is commutative. Number fields,
i.e. sets of numbers closed under addition, multiplica-
tion, subtraction and division by non-zero numbers,
were implicitly included in the very first studies of alge-
braic equations. Associative-commutative rings and
74
fields are the main objects studied in commutative alge-
bra and in the closely related field algebraic geometry.
Another important type of algebra with two binary
operations is a lattice. Typical examples of lattices
include: the system of subsets of a given set with the
operations of set-theoretic union and intersection, and
the set of positive integers with the operations of taking
the least common multiple and the greatest common
divisor.
Linear (or vector) spaces over a field may be treated
as universal algebras with one binary operation (addi-
tion) and with a unary operation (multiplication by
scalars of the ground field). Linear spaces over skew-
fields have also been studied. If a ring is considered
instead of a set of scalars, the more general concept of
a module is obtained. An important part of algebra,
linear algebra, studies linear spaces, modules and their
linear transformations, as well as problems related to
them. A part of it, the theory of linear equations and
the theory of matrices, was formulated as early as the
nineteenth century. A closely related subject is that of
multilinear algebra.
Initial studies on the general theory of arbitrary
universal algebras (this theory is sometimes called
'universal algebra') date back to the nineteen-thirties
and were carried out by G. Birkhoff. At the same time
A.I. Mal'tsev and A. Tarski laid the foundations for the
theory of models (cf. Model), i.e. sets with marked rela-
tions on them. Subsequently, the theory of universal
algebras and the theory of models became so closely
linked that they gave rise to a new discipline, inter-
mediate between algebra and mathematical logic, called
the theory of algebraic systems (cf. Algebraic system),
the subject of which are sets with algebraic operations
and relations defined on them.
A number of disciplines intermediate between alge-
bra and other fields of mathematics have been created
by the introduction into universal algebras of comple-
mentary structures compatible with the algebraic opera-
tions. These include topological algebra (including the
theory of topological groups and Lie groups, cf. Lie
group; Topological group), the theory of normed rings
(cf. Normed ring), differential algebra, and theories of
various ordered algebraic formations. Homological alge-
bra, which originates both from algebra and from
topology, arose in the nineteen-fifties as a discipline in
its own right.
The role of algebra in modern mathematics is
extremely important, and there is a trend towards
further 'algebraization' of mathematics. A typical way
of studying many mathematical objects that are some-
times far removed from algebra is to construct alge-
braic systems which adequately represent the behaviour
of these objects. Thus, the study of Lie groups can be

largely reduced to the study of their algebraic counter-
parts: Lie algebras (cf. Lie algebra). A similar method
is used in topology: to each topological space is
assigned, in some standard manner, an infinite series of
homology groups (cf. Homology group), and these series
of algebraic 'reflections of them' make it possible to
evaluate, very accurately, the properties of the spaces
themselves. The recent major discoveries in topology
were made using algebra as a tool (cf. Algebraic topol-
ogy).
It would appear at first sight that the translation of
problems into the language of algebra, solving them in
this language and translating them back is merely a
superfluous complication. In fact, such a method turns
out to be highly convenient, and occasionally the only
possible one. This is because by algebraization one
solves problems not only by purely verbal
considerations, but also by using the powerful
apparatus of formal algebraic calculations, so that one
may occasionally overcome highly involved complica-
tions. This role of algebra in mathematics may be com-
pared with the role of modern computers in the solu-
tion of practical problems.
Algebraic concepts and methods are widely employed
in number theory (cf. Algebraic number theory), in
functional analysis, in the theory of differential equa-
tions, in geometry (cf. Invariants, theory of; Projective
geometry; Tensor algebra), and in other mathematical
disciplines.
Besides its fundamental role in mathematics, algebra
is very important from the point of view of its applica-
tions; examples are applications within physics (the
representation theory of finite groups in quantum
mechanics; discrete groups in crystallography), cyber-
netics (cf. Automata, theory of) and mathematical
economics (linear inequalities, d. Linear inequality).
References
[I] The history of mathematics from the earliest times to the begin-
ning of the 19-th century, Vol. 1-3, Moscow, 1970-1972 (in Rus-
sian).
[2] MAL'TSEV, A.I.: 'On the history of algebra in the USSR during
her first twenty-five years', Algebra and LogiC 10, no. I (1971),
68-75. (Algebra i Logika 10, no. I (1971), 103-118)
[3] Mathematics, its content, methods and meaning, Vol. 1-3, Mos-
cow, 1956 (in Russian). Collection of articles.
[4] KUROSH, A.G.: Higher algebra, Mir, 1972 (in Russian).
[5] BOURBAKI, N.: Elements of mathematics. Algebra: AlgebraiC
structures. Linear algebra, I, Addison-Wesley, 1974,
Chapt. I; 2 (translated from the French).
[6] WAERDEN, B.L. VAN DER: Algebra, 1-2, Springer, 1967-1971
(translated from the German).
[7] LANG, S.: Algebra, Addison-Wesley, 1974.
[8] MAL'TSEv, A.I.: Algebraic systems, Springer, 1973 (translated
from the Russian).
For references, see also the articles on individual
algebraic disciplines.
Yu.I. Merzlyakov
A.I. Shirshov
ALGEBRA OF FUNCTIONS
Editorial comments.
References
[A1] LIDL, R. and PILZ, G.: Applied abstract algebra, Springer,
1984.
[A2] JACOBSON, N.: Lectures in abstract algebra, 1-3, v. Nostrand,
1951-1964.
AMS 1980 Subject Classification: 01-XX, 10-XX, 12-
XX, 13-XX, 15-XX, 20-XX
ALGEBRA, C * - - See C -algebra.
AMS 1980 Subject Classification: 22D25, 46L05
ALGEBRA, FUNDAMENTAL THEOREM OF - The
theorem that states that any polynomial with complex
coefficients has a root in the field of complex numbers.
The theorem was first stated by A. Girard in 1629 and
by R. Descartes in 1637 in a formulation different from
the one employed today. C. MacLaurin and L. Euler
made the formulation more precise and gave it a form
which is equivalent to the one in use today: Any poly-
nomial with real coefficients can be decomposed into a
product of linear and quadratic factors with real coeffi-
cients. A proof of the fundamental theorem of algebra
was first given by J. d' Alembert in 1746. There fol-
lowed the proofs of Euler, P. Laplace, lL. Lagrange
and others in the second half of the eighteenth century.
All these proofs are based on the assumption that some
'ideal' roots of the polynomial in fact exist, after which
it is demonstrated that at least one of them is a com-
plex number. C.F. Gauss was the first to prove the fun-
damental theorem of algebra without basing himself on
the assumption that the roots do in fact exist. His proof
essentially consists of constructing the splitting field of
a polynomial. All proofs of the theorem involve some
form of topological properties of real and complex
numbers. The role of topology has ultimately been
reduced to the single assumption that a polynomial of
odd degree with real coefficients has a real root.
References
[I] KUROSH, A.G.: Higher algebra, Mir, 1972 (translated from the
Russian).
[2] LANG, S.: Algebra, Addison-Wesley, 1974.
[3] BASHMAKOVA, I.G.: 'On a proof of the fundamental theorem of
algebra', Istor. Mat. Issled., no. 10 (1957), 257-304 (in Rus-
sian). VN
. . RemesIenm'k ov
Editorial comments. For a proof based on the Brouwer
fixed-point theorem cf. [A 1].
References
[A 1] ARNOLD, B.H.: 'A topological proof of the fundamental
theorem of algebra', Amer. Math. Monthly 56 (1949), 465-
466.
AMS 1980 Subject Classification: 26CXX, 12005
ALGEBRA OF FUNCflONS, function algebra - A
semi-simple commutative Banach algebra A, realized as
75
ALGEBRA OF FUNCTIONS
an algebra of continuous functions on the space of
maximal ideals 9n. If a EA and if / is some function
defined on the spectrum of the element a (i.e. on the set
of values of the function a=a), then/(a) is some func-
tion on 9n. Clearly, it is not necessarily true that
/(a)EA. If, however, / is an entire function, then
/(a)EA for any aEA. The use of the Cauchy integral
formula permits a considerable strengthening of this
result: If the function / is analytic in some neighbour-
hood of the spectrum of the element a, then /(a)EA
and the mapping / ~/ (a) is a homomorphism of the
algebra of functions which are analytic in some neigh-
bourhood of the spectrum of a EA into A. This proposi-
tion is valid for non-semi-simple commutative Banach
algebras as well. Moreover, this class of functions that
are analytic in a neighbourhood of the spectrum of a
given element cannot be enlarged, in general. For
example, if A =L1(Z) and /(a)EA for all aEA with
spectrum in the interval [0, 1], then / is analytic in
some neighbourhood of this interval.
In a few cases lea) can also be defined for multi-
valued analytic functions f, but such a definition has
inherent difficulties. Thus, let A be the algebra of con-
tinuous functions in the disc I Z I ~ 1 that are analytic
in the disc I z I < 1 and that satisfy the condition
/ (0) = O. The unit disc is naturally identified with the
space of maximal ideals of A. The function /I(Z)=Z,
which is continuous on the space of maximal ideals,
does not belong to A, but is a solution of the quadratic
equation
fl- Z 2 = 0,
where z2EA.
If A is a semi-simple algebra with space of maximal
ideals X, if /EC(X) and if
p(j) - j"+alj"-l + ... +an = 0, aiEA,
with p' (f) Ef.(A ), the group of units of A (a simple
root), then /EA. Similarly, if /EC(X) and if
exp(f)EA, then/EA.
A function algebra is said to be a uniformly-
convergent algebra (or uniform algebra) if the norm in
this algebra defines a notion of convergence equivalent
to the uniform convergence of the functions on the a
space of maximal ideals. If II a 2 II = II a 112 for all a EA,
then A is a uniform algebra. The general example of a
uniform algebra is a closed subalgebra of the algebra of
bounded continuous functions on some topological
space, provided with the natural sup-norm.
If A is a uniform algebra and if its space of maximal
ideals is metrizable, then among the boundaries (not
only the closed ones) there is a minimal boundary f 0,
the closure of which is the Shilov boundary. The set fo
consists of 'peak points': Xo is a peak point if there
exists a function /EA such that Ilex) I < I I(xo) I for
76
all x=#=xo. In the present case any point in the space of
maximal ideals has a representing measure concen-
trated on fo.
A function algebra is said to be analytic if all func-
tions of this algebra that vanish on a non-empty open
subset of the space of maximal ideals vanish identically.
Algebras that are analytic with respect to the boundary
are defined in a similar manner. Any analytic algebra is
analytic with respect to the Shilov boundary; the con-
verse is usually not true.
A function algebra A is said to be regular if, for any
closed set F in the space X of maximal ideals of A and
for any point Xo not contained in F, it is possible to
find a function / EA such that / (x) = 1 for all x EF and
/(xo)=O. All regular algebras are normal, i.e. for any
pair of non-intersecting closed sets F, F 0 EX there
exists an element/EA such that/(x)=1 for all XEF
and I (x) = 0 for all x EF o. In a regular algebra, for
any finite open covering {U;}, 1:S;;;i~m, of the space X
there exists a partition of unity belonging to A, i.e. a
system of functions /J, ... ,f" EA for which
fl(x)+ ... +f,,(x) 1
and
Ji(x) =0 if x f1- Ui.
A function g is said to belong locally to the function
algebra A if for any point Xo EX there exists a neigh-
bourhood in which this function coincides with some
function of the algebra. Any function which locally
belongs to a regular algebra is itself an element of this
algebra.
An element / of a function algebra is called real if
/(x) is real for all x EX. If A is an algebra with real
generators /a and if
d
f In II exp(itfa) II -t-2
TOO
-00 1 +t
< 00
for all la, then A is regular.
An ideal in a Banach algebra is said to be primary if
it is contained in only one maximal ideal. If A is a reg-
ular function algebra, then each maximal ideal Xo con-
tains a smallest closed primary ideal J(xo) which is
contained in any closed primary ideal contained in Xo.
The ideal J(xo) is the closure of the ideal formed by
the functions lEA that vanish in some neighbourhood
(depending on f) of x 0 EX.
In the algebra of absolutely convergent Fourier series
with an adjoined identity any maximal ideal coincides
with the corresponding primary ideal.
Let A be a closed subalgebra of the algebra C (X),
where X is a compactum (which does not necessarily
coincide with the space of maximal ideals of A). Let A
separate the points of X, i.e. for any two different
points Xl, X 2 E X there exists a function I in A for
which I(x d=#=I(x 2) The algebra A is called symmetric

if both a function I (x) and the function I(x) belong to
it. According to the Stone- Weierstrass theorem, if A is
symmetric, then A =C(X). The algebra A called anti-
symmetric if it follows from the conditions f,lEA that I
is a constant function. In particular, algebras of ana-
lytic functions are anti-symmetric. A subset S eX is
called a set 01 anti-symmetry (with respect to the alge-
bra A) if any function lEA that is real on S is constant
on this set. It follows from this definition that the alge-
bra A is anti-symmetric if the whole set X is a set of
anti-symmetry. In the general case the space X can be
represented as the union of non-intersecting, closed,
maximal sets of anti-symmetry. Each maximal set of
anti-symmetry is an intersection of peak sets (a set P is
called a peak set if there exists a function lEA such
that II p = 1 and II (x) I < 1 if x tl.P). It follows that
the restriction A I y of the algebra A to a maximal set
of anti-symmetry is a closed (anti-symmetric) subalge-
bra of the algebra C (Y). If X is the space of maximal
ideals of the algebra A, the maximal sets of anti-
symmetry are connected. If a continuous function is
such that on each maximal set of anti-symmetry it coin-
cides with some function in the algebra A, then the
function itself belongs to A. This generalization of the
Stone- Weierstrass theorem makes it possible, in prin-
ciple, to reduce the study of arbitrary uniform algebras
to the study of anti-symmetric algebras A. However,
the study of arbitrary algebras A cannot be reduced to
the study of analytic algebras: There exists an example
of an algebra of type R (X) (a closed subalgebra of the
algebra C(X which does not coincide with C(X), and
is anti-symmetric and regular.
Let Re(A) be the real space of functions of the form
Re(f), where lEA; if Re(A) is an algebra or if Re(A) is
closed in C(X), then A =C(X). The space X can be
regarded as a part of the space of maximal ideals of the
algebra A; accordingly, not only the ordinary topology
of the space of maximal ideals, but also the metric
induced by the imbedding of X into the dual space to
A can be considered on X. The distance in the sense of
this metric will be denoted by PA. For any points
x I> X2 EX the inequality PA (x I> x2):S;;;2 is valid; the
relation PA(XI> x2)<2 is an equivalence relation, and
the equivalence classes are known as Gleason parts. If X
is the disc I z I :s;;; 1 and A is the closed sub algebra in
C (X) consisting of the functions analytic in I z I < 1,
then the metric PA is non-Euclidean, and the one-point
sets on the circle and in the interior of the disc serve as
the Gleason parts. Gleason parts do not necessarily
have an analytic structure: Any a-compact completely-
regular space is homeomorphic to the Gleason part of
the space of maximal ideals of some algebra, such that
the restriction of the algebra to this part contains all
bounded continuous functions. The fact that two points
ALGEBRA OF LOGIC
belong to the same Gleason part can be described in
terms of the representing measures on the Shilov boun-
dary: Two such points have two mutually absolutely
continuous representing measures with bounded deriva-
tives. An algebra for which Re(A Ir) is dense in C(f)
is called a Dirichlet algebra; if P is a Gleason part in
the space of maximal ideals of a Dirichlet algebra
which contains more than one point, then there exists a
continuous one-to-one mapping 1/1 of the disc I z I < I
into P such that for any function lEA the function
I(I/I(z is analytic in I z I <l. Thus, P has a structure
with respect to which the functions lEA are analytic;
the mapping 1/1 is not a homeomorphism in general if P
is endowed with the ordinary topology of the space of
maximal ideals, but 1/1 is a homeomorphism if P is
endowed with the metric PA.
For references see Banach algebra.
E.A. Gorin
Editorial comments.
References
[A1] GAMELIN, T.W.: Uniform algebras, Prentice-Hall, 1969.
[A2] STOUT, E.L.: The theory of uniform algebras, Bogden and
Quigley, 1971.
AMS 1980 Subject Classification: 46J 10
ALGEBRA OF LOGIC - The branch of mathematical
logic that deals with propositions from the aspect of
their logical meanings (true or false) and with logical
operations on them.
The algebra of logic originated in the middle of the
nineteenth century with the studies of G. Boole [1], [2],
and was subsequently developed by C.S. Peirce, P.S.
Poretskii, B. Russell, D. Hilbert, and others. The
development of the algebra of logic was an attempt to
solve traditional logical problems by algebraic methods.
With the birth of the theory of sets (in the eighteen-
seventies), propositions and logical operations on them
became the principal subject of algebra of logic. Propo-
sitions are understood to mean statements about which
it is meaningful to ask whether they are true or false.
For instance, the proposition 'a whale is an animal' is
true, while the statement 'all angles are right angles' is
false. The connectives 'and', 'or', 'if ... then', 'is
equivalent to', the particle 'not', etc., which are com-
monly used in the language of logic, make it possible to
construct new, more 'complicated', propositions from
those already established. Thus, given that 'x> 2' and
'x:S;;;3', it is possible to obtain, by using the connective
'and', the proposition 'x>2 and x:S;;;3'; by using the
connective 'or' it is possible to obtain the proposition
'x>2 or x:S;;;3', etc.
The truth or the falsehood of the propositions
obtained in this way depend on the truth or falsehood
of the initial propositions and on a corresponding treat-
ment of the connectives as operations on propositions.
77
ALGEBRA OF LOGIC

Often truth is symbolized by the digit '1' and falsehood below:

by the digit '0'. The connectives 'and', 'or', 'if '" x y x x&y xVy x~y x-y
then', and 'is equivalent to' are denoted, respectively,
0 0 1 0 0 1 1
by the symbols & (conjunction), V (disjunction), ~
0 1 1 0 I 1 0
(implication), and (equivalence); negation is
I 0 0 0 I 0 0
represented by a bar - above the symbol. In addition to
1 1 0 1 1 1 1
individual propositions, examples of which are given
above, one also uses variable propositions, i.e. variables Tables for arbitrary functions of the algebra of logic
whose values may be any individual statements given in are constructed in a similar manner. This is the so-
advance. As the next step, the concept of a formula called tabular way of specifying functions of the algebra
was inductively introduced. It formalizes the concept of of logic. The tables themselves are sometimes called
a compound proposition. Let A, B, C, ... , denote truth tables. The following equalities play an important
individual propositions, and let x, y, z, ... , denote part in the transformations of formulas into equivalent
variable propositions. Each one of the above letters is formulas:
known as a formula. Let the symbol * denote anyone x&y = y&x, xVy = yVx (1)
of the connectives listed above, and let ~ and ?E denote (the law of commutativity);
formulas; then (~*?E) and i will be formulas (e.g. (x&y)&z = x&(y&z), (xVy)Vz = xV(yVz) (2)
x&y)~Z). The connectives and the particle 'not' are
(the law of associativity);
looked upon as operations on quantities which assume
the values 0 and 1, the results of these operations are x&(xVy) = x, xV(x&y) = x (3)
also the digits 0 and 1. The conjunction x&y is equal to (the law of absorption);
1 if and only if both x and y equal 1; the disjunction x&(yVz) = (x&y)V(x&z), (4)
xVy equals 0 if and only if x and y are both equal to
xV(y&Z) = (xVy)&(xVz)
0; the implication x~y is equal to 0 if and only if x = I
and y = 0; the equivalence x -y equals 1 if and only if (the laws of distributivity);
the values of x and y are identical; and the negation X- x&x = 0 (5)
is equal to 1 if and only if x = O. If the values of the
(the law of contradiction);
propositions in each formula are given, the value 0 or 1
can be assigned to the formula. This means that any xVx = 1 (6)
formula can be considered as a way of stating or as a (the law of the excluded middle);
way of realizing a function of the algebra of logic, i.e. a
x~y = xVy, (7)
function which is defined on samples of zeros and ones,
and with as values also 0 or 1. Two formulas ~ and ?E x~y = (x&y) V (x&y).

are said to be equal (~=?E) if they realize equal func- If these equalities are used, it is possible to obtain new
tions. The subject matter of the algebra of logic is the equalities, even without the use of tables. These new
treatment of functions of the algebra of logic and the equalities are ottained by the so-called identity transfor-
operations on these functions. In what follows below, mations which, generally speaking, alter the expression,
the class of functions of the algebra of logic will be but not the function realized by the expression. For
extended to the class of functions whose arguments, as instance, the absorption laws yield the law of idempo-
well as the functions themselves, assume the values of tency x V x =x. The above equalities frequently sim-
elements of a finite fixed set E; the number of opera- plify the notation of formulas by eliminating some of
tions on the functions will also be extended. The alge- the parentheses. Thus, the relations (1) and (2) make it
bra of logic is sometimes understood to mean this latter possible to replace the formulas
concept. However, from the aspect of practical applica- ... (~l &~2)& ... )&~s) and
tions, the most important case is that of sets E of cardi- ... (?E 1 V?E 2 )V ... )V?Es) by the more compact
nality two, and this case will therefore be discussed in notation ~1&~2&'" &~s and ?El V?E2 V ... V?Es'
greatest detail. The results presented here are also The former relation is known as the conjunction of the
closely connected with a second approach to the study factors ~t. ... ,~s, while the latter is called the disjunc-
of propositions - the so-called propositional calculus. tion of the terms ?E 1 , ,?Es' It is also seen from the
In specifying the functions of the algebra of logic one equalities (5), (6) and (7) that, if the constants 0 and I,
often uses tables which contain all the combinations of the implication and the equivalence, are regarded as
the values of the variables and the values of the func- functions, they may be expressed by conjunctions, dis-
tions on these combinations. Thus, the synoptic table of junctions and negations. Thus, any function of the alge-
the functions X-, x&y, xVy, x~y, and x-y is given bra of logic may be realized by a formula written using

78

only the symbols &, V and-.
The set of all formulas involving variable proposi-
tions, some of the symbols &, V, ~, "', -, and the
constants 0 and 1 is called a language over these symbols
and constants. Equations (1) - (7) show that, for any
formula in the language over &, V, ~, "', -, 0, and 1,
it is possible to find an equivalent formula in the
language over &, V, -, 0, and 1; for example
(x~y)~z = x Vy)&Z)Vx Vy)&z).
In the latter language a special role is played by a class
of formulas which may be written in the form
2r\ V ... V2r s, 0 or 1, where s;;;;.l and where each 2rj
is either a variable proposition, its negation or a con-
junction of them, so that none of the 2rj contains equal
x
factors or factors of the form x and together, and all
2rj are pairwise unequal. Here the parentheses are omit-
ted, since it is assumed that the conjunction operation
is 'stronger' than disjunction, i.e. when calculating with
respect to the given values of the variables, the values
of 2r \, ... , 2rs are calculated first. These expressions
are called disjunctive normal forms. Each formula u in a
language over &, V, ~, "', -, 0, 1 which realizes a
non-zero function of the algebra of logic, may be con-
verted with the aid of the equalities (1) - (7) into an
equivalent disjunctive normal form, containirJg all vari-
able formulas of u and any number of other variables,
and each 2rj in this disjunctive normal form containirJg
the same variables. Such a disjunctive normal form is
said to be a perfect disjunctive normal form of the for-
mula u; the perfect disjunctive normal form of 0 is the
formula 0 itself. The possibility of reduction to a per-
fect disjunctive normal form forms the basis of an algo-
rithm determining the equality or inequality of two
given formulas. The procedure of this algorithm is as
follows: The given formulas u \ and U2 are reduced to
perfect disjunctive normal forms which contain all the
variables comprised in both u \ and U2 and the two
expressions are compared; if they are identical, u\ =U2;
if not, u \ =;FU 2. An important role in the algebra of
logic and its applications is played by contracted dis-
junctive normal forms, i.e. disjunctive normal forms
which satisfy the following conditions: 1) it does not
contain pairs of elements 2rj and 2rj such that any fac-
tor of 2rj is contained in 2rj ; 2) for any two elements 2rj
and 2rj one of which contains some variable as factor,
while the other contains the negation of this variable
(under the condition that the pair of elements does not
contain a second variable for which the same is true),
there exists (in the same disjunctive normal form) an
element 2rk which is equal to the conjunction of the
other factors of these two elements. Any disjunctive
normal form may be reduced with the aid of the equali-
ties (1) - (7) to an equal contracted disjunctive normal
ALGEBRA OF LOGIC
form. Thus, the contracted disjunctive normal form for
the formula x"'(y~z~(x&z is x&yVzVx&y. The
formulas u \ and U2 are equal if and only if their con-
tracted disjunctive normal forms are identical. In addi-
tion to the disjunctive normal form, use is also made of
the conjunctive normal form, i.e. of expressions which
are obtained from a disjunctive normal form by substi-
tuting the symbol V for &, & for V, and 0 for 1. For
instance, the disjunctive normal form x&y VX&Z yields
the conjunctive normal form (x Vy)&(xV z). An opera-
tion (or function) f is said to be- the dual of an opera-
tion I/; if the table defining f is obtained from the table
which defines I/; by substituting 1 for 0 and 0 for 1
throughout (including interchange in the values of the
functions). Thus, the conjunction and the disjunction
are mutually dual, the negation is dual to itself, the
constants 0 and 1 are mutually dual, etc. A transforma-
tion of a formula which consists in replacing the sym-
bols of all operations and expressions by the symbols
of their dual operations, and in exchanging 0 for 1 and
1 for 0, is called a duality transformation. If the equality
2r = ~ is true, and if 2r. is the dual of 2r while ~. is the
dual of ~, then it is also true that 2r. =~., and this
equality is called the dual of the above equality; this is
the so-called duality prinCiple. The pairs of the laws (1),
(2) and (3) are examples of dual equalities; equality (5)
is dual to equality (6), and each conjunctive normal
form is the dual of some disjunctive normal form. A
perfect conjunctive normal form and a contracted con-
junctive normal form are defined as conjunctive normal
forms such that their dual expressions are a perfect dis-
junctive normal form and a contracted disjunctive nor-
mal form, respectively. Perfect and contracted disjunc-
tive and conjunctive normal forms are used for survey-
ing all hypotheses and all consequences of a given for-
mula. A hypothesis of a formula 2r is understood to be a
formula ~ such that (~~2r) = I; a consequence of a for-
mula 2r is a formula ~ such that (2r~~)= 1. A
hypothesis of a formula 2r is said to be simple if it is a
conjunction of variable formulas or their negations and
if, after discarding anyone of its factors, it is no longer
a hypothesis of the formula 2r. Similarly, a consequence
of 2r is called simple if it is a disjunction of variable
formulas or their negations and if, after one of its ele-
ments is discarded, it is no longer a consequence of 2r.
The survey of hypotheses and consequences is based on
the indication of an algorithm which transforms a given
formula into all of its simple hypotheses and
consequences and, with the aid of the laws (2) - (7),
into all remainirJg hypotheses and consequences. The
algorithm is based on the following facts. If 2r =~, then
2r and ~ have identical hypotheses and identical conse-
quences. An element of a disjunctive normal form is a
hypothesis of that disjunctive normal form, while a fac-
79
ALGEBRA OF LOGIC
tor of a conjunctive normal form is a consequence of
that conjunctive normal form. If m: is a hypothesis of a
proposition 58, then m:&~ is also a hypothesis for 58; if
m: is a consequence of 58, then m:V~ is also a conse-
quence of 58. If m: and ~ are hypotheses of a proposi-
tion 58, then m:V ~ is also a hypothesis for 58; if m: and
~ are consequences of 58, then m:& ~ is also a conse-
quence of 58. A perfect disjunctive normal form has no
hypotheses (not containing letters which are not con-
tained in that disjunctive normal form) other than the
disjunctions of some of its elements or of disjunctive
normal forms equal to them. A perfect conjunctive nor-
mal form has no consequences other than the conjunc-
tions of some of its factors or of propositions equal to
them. A contracted disjunctive normal form is the dis-
junction of all its simple hypotheses; a contracted con-
junctive normal form is the conjunction of all its simple
consequences. A contracted disjunctive normal form
has important applications. The first problem which
should be mentioned is the minimization of functions
of the algebra of logic, which is a part of the problem
of synthesis of control (switching) systems. The minimi-
zation offunctions of the algebra of logic consists in con-
structing a disjunctive normal form for the given func-
tion of the algebra of logic which realizes this function
and has the smallest sum of the number of factors in its
elements, i.e. has minimal complexity. Such disjunctive
normal forms are called minimal. Each minimal dis-
junctive normal form for a given function of the alge-
bra of logic which is not a constant is obtained from
the contracted disjunctive normal form of this function
by discarding a number of elements. For certain func-
tions the contracted disjunctive normal form may coin-
cide with the minimal disjunctive normal form. Such a
case occurs, for example, with monotone functions, i.e.
functions realizable by formulas over &, V, 0, 1.
In the language over &, V, ~, ~, 0, 1, +, when the
sign + is interpreted as addition modulo 2, the follow-
ing relations are valid:
xVy = x&y)+x)+y, (8)
x~y = i&y, x~y = (x+y)+I, (9)
x+y = (x&y) V (i&y), I = xVi. (10)
These formulas make it possible to translate formulas
from the language over &, V, ~, -, 0, and 1 into
equivalent formulas in the language over &, +, 1, and
vice versa. The identity transformations in the latter
language are realized with the aid of equalities esta-
blished for the conjunction and the following additional
equalities:
x+y = y+x, (II)
(x+y)+z = x+(y+z), (12)
x&(y+z) = x&y+x&z, (13)
x&y = x, x+(y+y) = x, x&1 = x. (14)
80
Here it is considered, as before, that conjunction is a
stronger connective than the sign +. These equalities
are sufficient to deduce any valid equality in the
language over &, +, 1, with the aid of identity
transformations, just as in the case of the language over
&, V, ~, "', -, +, 1. A proposition in this language is
called a reduced polynomial if it is of the form
m: 1 + ... +m:s , when m: j is equal to 1 or is a variable or
a conjunction of various variables without negations,
m:j=#=m:j if i=/=j, s;;:.l, or else is equal to 1+1. Thus, the
expression x&y&z +x&y + 1 is a reduced polynomial.
Any formula of the algebra of logic may be converted
to a reduced polynomial by identity transformations.
The equality m: = 58 is valid if and only if the reduced
polynomial for m: is identical with the reduced polyno-
mial for 58. In addition to the languages mentioned
above there are also other languages equivalent to
them. Two languages are called equivalent if each for-
mula in one language can be converted, by means of
certain conversion rules, into another equivalent for-
mula in the second language, and vice versa. It is suffi-
cient to base such a language on an arbitrary system of
operations (and constants) such that the operations
(and constants) of the system may be used to represent
any function of the algebra of logic. Such systems are
said to be functionally complete. Examples of complete
systems are {xVy}, {xVy,x}, {x+y, l,x&y}, etc.
There exists an algorithm which can be used to estab-
lish the completeness or incompleteness of an arbitrary
finite system of functions of the algebra of logic. It is
based on the following fact. A system of functions of
the algebra of logic is complete if and only if it con-
tains functions fI (x, y, ... , v) and 12 (x, y, ... , v) such
that fl (0,0, ... ,0)= 1 and 12(1, 1, ... ,1)=0, as well
as functions h, f4 and f5, where h =#=h' f4 is not
monotone, and the reduced polynomial of f5 contains
an element m: with more than one factor. There are
also other languages based on systems of operations
which are not functionally complete, and the number of
such languages is infinite. They contain infinitely many
pairwise incomparable languages (in the sense that
there is no way of translating from one language to
another by means of identity transformations). How-
ever, for any language based on some operations of the
algebra of logic there exists a finite system of equations
in this language such that any equality can be deduced
with the aid of identity transformations of the equa-
tions of this system. Such a system called a deductively
complete system of equalities in this language. Thus, the
equalities (1) - (6) constitute a complete system of
equalities of the language over &, V, -, 0, 1.
In considering one of the languages discussed above
in conjunction with some complete system of equalities
in this language, a tabular statement of the basic opera-
tions in the language and the requirement that proposi-

tions be the values of its variables are sometimes dis-
carded. Instead, various interpretations of the language
are permitted. These interpretations involve some set of
objects (which are used as values of the variables) and
some system of operations over the objects of this set,
which satisfy the equalities from a complete system of
equalities of the language. Thus, such a procedure
results in the transformation of the language over &,
V, -, 0, I into the language of a Boolean algebra; the
language over &, +, I is transformed into the language
of a Boolean ring (with a unit element); the language
over &, V, - is transformed into the language of a dis-
tributive lattice, etc.
Historically, the development of the algebra of logic
has been stimulated mainly by the problems to which it
has been applied. An important application is the
theory of electric circuits. In certain cases circuits can-
not be described in terms of the ordinary two-valued
algebra of logic, and a many-valued logic must be con-
sidered (cf. Many-valued logic).
References
[I) BOOLE, G.: The mathematical analysis of logic: being an esstry
towards a calculus of deductive reasoning, Macmillan, 1847.
[2) BOOLE, G.: An investigation of the laws of thought, on which are
founded the mathematical theories of logic and probabilities,
Dover, reprint, 1951.
[3) PORETSKII, P.S.: 'On methods of solving logical equalities and
on an inverse method of mathematical logic', in Collection of
protocols of sessions of the FizMat. section of the Society of
Natural Scientists at Kazan. Univ., Vol. 2, Kazan', 1884, pp.
161-330 (in Russian).
[4) HILBERT, D. and ACKERMAN, W.: Grundzuge der theoretischen
Logik, Dover, reprint, 1946.
[5) NOVIKOV, P.S.: Elements of mathematical logic, Edinburgh,
1964 (translated from the Russian).
[6) YABLONSKIi, S.V., GAVRILOV, G.P. and KUDRYAVTSEV, V.B.:
Functions of the algebra of logic and Post classes, Moscow, 1964
(in Russian).
V.B. Kudryavtsev
Editorial comments. In the theory of Boolean functions
contracted disjunctive normal forms are also called abridged
disjunctive normal forms, cf. Boolean functions, normal
forms of.
AMS 1980 Subject Classification: 03GXX
ALGEBRA OF MEASURES, measure algebra - The
algebra M (G) of complex-valued regular Borel meas-
ures on a locally compact Abelian group G that have
bounded variation, with the ordinary linear operations
and with convolution as multiplication (cf. Harmonic
analysis, abstract). The convolution A*J.L of the meas-
ures A, J.LEM(G) is completely defined by the condition
that, for any continuous function f on G with compact
support,
jfd(h*P.) = j jf(x+y)dh(X)dp.(y).
G GG
If the total variation of a measure is taken as norm,
M (G) becomes a commutative Banach algebra over the
field of complex numbers. The algebra of measures
ALGEBRA OF MEASURES
M (G) has a unit which is the 8-measure located at the
zero of the group. The set of discrete measures con-
tained in M (G) forms a closed subalgebra.
To each function f which belongs to the group alge-
bra L 1(G) may be assigned a corresponding measure
J.LfEM(G) in accordance with the rule
p./E) = ! dx
(integral with respect to the Haar measure). The result
is an isometric isomorphic imbedding Ll(G)~M(G).
Under this imbedding the image is a closed ideal in
M(G).
The Fourier- Stieltjes transform of a measure
J.LEM(G) is the function ~ on the dual group Gdefined
by the formula
Then ~=~.~ and II J.L II =0 if ~ O. In particular,
M (G) is an algebra without a radical.
If the group G is not discrete, then the structure of
the measure algebra M(G) is extremely complicated: It
is not symmetric and its space of maximal ideals has a
number of pathological properties. For instance, this
space contains infinite-dimens~ona1 analytic sets, and
the naturally imbedded group G in it is not dense even
in the Shilov boundary. Nevertheless, the idempotent
measures, i.e. measures for which J.L*J.L=J.L, are known.
Each idempotent measure is a finite integer combina-
tion n 1J.Ll + ... + nkJ.Lb where J.Li = Xi Pi, and where Pi is
the Haar measure of a compact subgroup, and Xi is a
character. In the case G = Z this means that a sequence
(cm ) of zeros and ones is the Fourier- Stieltjes
transform of some measure on the circle if and only if
(cm ) differs from a periodic sequence by not more than
a finite number of terms.
In the general case the theorem on idempotent meas-
ures can be naturally interpreted in terms of the coho-
mology spaces of dimension zero of the space of maxi-
mal ideals. A satisfactory description is also known for
other cohomology groups of the space of maximal
ideals of the measure algebra. This makes it possible to
tell, in particular, if a logarithm of an invertible meas-
ure from M (G) can be taken (one-dimensional integral
cohomology).
References
[1) RUDIN, W.: Fourier analysis on groups, Interscience, 1962.
[2) TAYLOR, J.L.: 'The cohomology of the spectrum of a measure
algebra', Acta Math. 126 (1971), 195-225.
E.A. Gorin
Editorial comments.
References
[A 1] GRAHAM, c.c. and MCGEHEE, O.c.: Essays in commutative
harmonic analysis, Springer, 1979,
[A2] TAYLOR, J.L.: Measure algebras, Amer. Math. Soc., 1972.
AMS 1980 Subject Classification: 28A60, 43A20
81
ALGEBRA OF SETS
ALGEBRA OF SETS - A non-empty collection of
subsets of some set n that is closed under the set-
theoretic operations (of union, intersection, taking com-
plements), carried out a finite number of times. In
order for a class of subsets of a set n to be an algebra
of sets, it is necessary (and sufficient) for it to be closed
under finite unions and taking the complement. An
algebra of sets that is closed under countable unions is
known as a a-algebra of sets. Any a-algebra of sets is
closed under the set-theoretic operations carried out a
countable number of times.
Examples. 1) The collection of finite subsets of an
arbitrary set n and their complements is an algebra of
sets; the collection consisting of the at most countable
subsets of n and their complements is a a-algebra of
sets.
2) The collection of finite unions of intervals of the
type
forms an algebra of sets.
3) let n be a topological space; the a-algebra B of
sets generated by the open subsets of n (in other words,
the smallest a-algebra of sets containing all open sub-
sets of n) is known as the Borel a-algebra of subsets of
n, while the sets belonging to B are known as Borel
sets.
4) let n=RT , where T is an arbitrary set (i.e. n is the
set of all real functions on 1); the class A of sets of the
type
where E is a Borel subset of Rk, is an algebra of sets;
in the theory of random processes a probability measure
is often originally defined only on an algebra of this
type, and is their subsequently extended to a wider
class of sets (to the a-algebra generated by A).
5) The collection of Lebesgue-measurable subsets of
Rk is a a-algebra of sets.
Algebras (respectively, a-algebras) are the natural
domain of definition of finitely-additive (respectively,
a-additive) measures. According to the theorem of
extension of measures, any a-finite, a-additive measure,
defined on an algebra A, can be uniquely extended to a
a-additive measure defined on the a-algebra generated
by A.
References
[I] DUNFORD, N. and SCHWARTZ, J.T.: Linear operators. General
theory, I, Interscience, 1958.
[2] HALMOS, P.R.: Measure theory, v. Nostrand, 1950.
[3] NEVEU, J.: Bases mathematiques du calcul des probabilites, Mas-
son, 1970.
v. V. Sazonov
AMS 1980 Subject Classification: 28A05
ALGEBRA WITH ASSOCIATIVE POWERS - A
linear algebra A over a field F each element of which
82
generates an associative sub algebra. The set of all alge-
bras with associative powers over a given field F forms
a variety of algebras which, if the characteristic of the
field F is zero, is defined by the system of identities
(x, x, x) = (x 2 , x, x) = 0,
where (x,y, z)=(xy)z-x(yz). If Fis an infinite field of
prime characteristic p, then the variety of algebras with
associative powers cannot be defined by any finite sys-
tem of identities, but an independent, infinite system of
identities which defines it is known [3]. If a commuta-
tive algebra A with associative powers of characteristic
other than 2 has an idempotent e=;FO, then A can be
decomposed according to Peirce into a direct sum of vec-
tor subspaces:
A = Ao(e)EflA 1/2(e)ElM I(e), (*)
where A >..(e) = {a EA: ea =;\a}, ;\=0, 1/2, 1. Here
A o(e) and A I (e) are sub algebras, A o(e)A I (e)=O,
A 1/2(e)A 1/2(e)CAO(e)+A I(e),
A>..(e)AI/2(e)CA I / 2(e)+A j ->..(e) for ;\=0,1. The
decomposition (*) plays a fundamental role in the
structure theory of algebras with associative powers.
References
[I] ALBERT,A.A.: 'Power-associative rings', Trans. A mer. Math.
Soc. 64 (1948), 552-593.
[2] GAiNOV, A.T.: 'Identity relations for binary Lie rings', Uspekhi
Mat. Nauk 12, no. 3 (1957), 141-146 (in Russian).
[3] GAiNOV, A.T.: 'Power-associative algebras over a finite-
characteristic field', Algebra and Logic 9, no. 1 (1970),5-19.
(Algebra i Logika 9, no. 1 (1970), 9-33)
A.T. Gainov
Editorial comments. An algebra with associative powers
is also called a power-associative algebra. The fact that the
set of algebras with associative powers over a field of non-
zero characteristic forms a variety defined by
(x, x, x)=(x 2 , x, x)=O was proved in [A1].
References
[A 1] ALBERT, A.A.: 'On the power associativity of rings', Summa
Brasiliensis Math. 2 (1948), 21-33.
AMS 1980 Subject Classification: 17A 10, 17A05, 16-
XX
ALGEBRAIC ALGEBRA - An algebra A with associ-
ative powers (in particular, an associative algebra) over
a field F all elements of which are algebraic (an ele-
ment a EA is called algebraic if the subalgebra F[a]
generated by it is finite-dimensional or, equivalently, if
the element a has an annihilating polynomial with coef-
ficients from the ground field F). An algebra A is called
an algebraic algebra of bounded degree if it is algebraic
and if the set of degrees of the minimal annihilating
polynomials of its elements is bounded. Subalgebras
and homomorphic images of an algebraic algebra (of
bounded degree) are algebraic algebras (of bounded
degree).

Examples: locally finite algebras (in particular,
finite-dimensional ones), nil-algebras and associative
skew-fields with a countable set of generators over an
uncountable field.
The algebras considered below are associative. The
Jacobson radical of an algebraic algebra is a nil-ideal. A
primitive algebraic algebra A is isomorphic to a dense
algebra of linear transformations of a vector space over
a skew-field; if, in addition, A is of bounded degree,
then A is isomorphic to a ring of matrices over a skew-
field. An algebraic algebra without non-zero nilpotent
elements (in particular, a skew-field) over a finite field
is commutative. It follows that finite skew-fields are
commutative. An algebraic algebra of bounded degree
satisfies a polynomial identity (cf. PI-algebra). An alge-
braic PI-algebra is locally finite. If the ground field is
uncountable, then the algebras obtained from an alge-
braic algebra by extension of the ground field, and the
tensor product of algebraic algebras, are algebraic alge-
bras.
References
[I] JACOBSON, N.: Structure of rings, Arner. Math. Soc., 1956.
[2] HERSTEIN, I.: Noncommutative rings, Math. Assoc. Amer., 1968.
VN. Latyshev
AMS 1980 Subject Classification: 16A16, 12F05
ALGEBRAIC BRANCH POINT, algebraic singular
point - An isolated branch point a of finite order of an
analytic function 1 (z), having the property that the
limit limZ-+ 1 (z) exists for any regular element of con-
Q dimension one. An algebraic curve is the most fre-
tinuation of 1 in a domain for which a is a boundary
point. More exactly, a singular point a in the complex
z-plane for the complete analytic function 1 (z), under
continuation of some regular element eo of this func-
tion with centre z 0 along paths passing through a, is
called an algebraic branch point if it fulfills the follow-
ing conditions: I) There exists a positive number p such
that the element eo may be extended along an arbitrary
continuous curve lying in the annulus D {z: = tions of an irreducible algebraic curve over k is the
0< I z-a I <p}; 2) there exists a positive integer k>1
such that if z 1 is an arbitrary point of D, the analytic
continuation of the element eo in D yields exactly k
different elements of the function 1 (z) with centre z 1 ;
if e 1 is an arbitrary element with centre z 1, all the
remaining k - I elements with centre z 1 can be
obtained by analytic continuation along closed paths
around the point a; and 3) the values at the points z of
D of all elements which are obtainable from eo by con-
tinuation in D tend to a definite, finite or infinite, limit
as z tends to a while remaining in D.
The number k - I is said to be the order of the alge-
braic branch point. All branches of the function 1 (z)
obtainable by analytic continuation of the element eo
in the annulus D may be represented in a deleted
ALGEBRAIC CURVE
neighbourhood of a by a generalized Laurent series
(Puiseux series):
00
j(z) = ~ cn(z-at/ k , m~O.
n=-m
The point at infinity, a = 00, is called an algebraic
branch point for a function 1 (z) if the point b 0 is an =
algebraic branch point of the function g(w)= l(l / w).
There may exist several (and even an infinite number
of) different algebraic branch points and regular points
of a complete analytic function with a given affix a.
References
[1] MARKUSHEVICH, A.I.: Theory offunctions of a complex variable,
3, Chelsea, 1977, Chapt.8 (translated from the Russian).
[2] HURWITZ, A. and COURANT, R.: Vorlesungen iiber allgemeine
Funktionentheorie und elliptische Funktionen, 4, Springer, 1968.
E.D. Solomentsev
AMS 1980 Subject Classification: 30AXX, 32830
ALGEBRAIC CLOSURE 01 a field k - An algebraic
extension (cf. Extension of a field) of k that is an alge-
braically closed field. Such an extension exists for each
field and is unique up to isomorphism. The algebraic
closure of the field of real numbers is the field of com-
plex numbers (cf. Algebra, fundamental theorem of).
VN. Remeslennikov
AMS 1980 Subject Classification: 12F05
ALGEBRAIC CURVE - An algebraic variety of
quently studied object in algebraic geometry. In the
sequel, an algebraic curve means an irreducible alge-
braic curve over an algebraically closed field.
The simplest and clearest concept is that of a plane
affine algebraic curve. This is a set of points in an
affine plane A~ satisfying the equation l(x,y)=O,
where I(x, y) is a polynomial with coefficients from an
algebraically closed field k. The field of rational func-
field of algebraic functions in one variable of the form
k(x, y), where x and y are connected by the equation
l(x,y)=O, where I(x,y) is a polynomial over k. This
means that every algebraic curve is birationally iso-
morphic to a plane affine curve.
It has been known for a long time that even when
studying affine curves fundamental relationships can
only be revealed by considering points at an infinite
distance, and by a detailed study of the singular points.
In order to study all the points of an affine curve, the
curve is imbedded into a projective space pn, with sub-
seqlloent closure in the Zariski topology. In this way a
projective curve X is obtained, and the initial affine
curve Y may be obtained from X by deleting a finite
number of points. If Y is irreducible, then X and Y are
83
ALGEBRAIC CURVE
birationally isomorphic. Every complete algebraic curve
is projective. If X is a smooth projective curve, then all
valuation rings of the field k(x) are given by the local
rings (!) X> X EX (cf. Local ring). If two smooth projec-
tive curves are birationally equivalent, then they are
isomorphic. A normal algebraic curve is smooth. In
particular, any irreducible algebraic curve is biratio-
nally equivalent to a smooth projective curve. The pro-
jective model of an algebraic curve, obtained by the
process of normalization, lies in some space pn. Any
smooth projective curve is isomorphic to a curve
situated in p3. Any plane algebraic curve can be con-
verted by a Cremona transformation into a curve with
ordinary singular points.
Divisors on a smooth algebraic curve are represented
by linear combinations of points with integer coeffi-
cients
where nx=O for almost all x (cf. Divisor). If all nx;;'O,
the divisor D is called positive, or effective, which is
written as D ;;.0. The degree of the divisor D is the
number
degD = ~nx'
The principal divisors form a subgroup P(X) of the
group Div X of all divisors on X. The quotient group
Div(X) / P(X) is called the group of'divisor classes and
is denoted by CI(X). The group CI(X) is isomorphic to
the group Pic(X) of classes of one-dimensional vector
bundles on X (cf. Vector bundle, algebraic). The degree
of the principal divisors on a smooth projective curve is
zero, and thus all divisors in one class are of the same
degree. In particular, one can speak of the degree of a
divisor class, and of the subgroup Clo(X) of divisor
classes of degree zero. The following equality is valid:
Cl(X) / Clo(X) = Z.
For the line pI, CI(PI)=Z, ClO(pl)=O, i.e. any divisor
of degree zero is a principal divisor. This property is
characteristic for rational smooth projective curves.
For any complete algebraic curve X, the number
'7T=dimHI(X, (!) x) is known as the arithmetic genus of
the algebraic curve X. If X is smooth, '7T is identical with
the dimension of the space HO(X, Q~) of all regular dif-
ferential forms on X; this dimension is known as the
genus of X. By definition, the genus of an algebraic
curve is equal to the genus of its non-singular model.
For any non-negative integer g there exists an algebraic
curve of genus g. Rational curves are distinguished by
the equality g = O. If X is a projective plane curve of
order m, then
17= (m-l)(m-2)
2 '
and its genus is given by the formula:
g = (m-l)(m-2)_d
2 '
84
where d is a non-negative integer which is a measure of
the smoothness on X. If X has only ordinary double
points, d is simply the number of singular points.
In particular, the genus of a plane smooth projective
curve is given by
_ (m -l)(m -2)
g - 2 '
which means that not every smooth projective curve is
plane. For a curve X in space the following estimate
holds:
(n - 2)2
for even n,
4
(n -l)(n -3) (I)
for odd n,
4
where n is the degree of X. Curves of degree n of maxi-
mal genus exist for each value of n and lie on a quadric
(G. Halphen, 1870).
The degree of the canonical class Kx of a smooth
projective curve X is connected with the genus of the
curve by means of the formula degKx =2g-2. If a
smooth projective curve X lies on a smooth algebraic
surface F, the adjunction formula Kx=X(X+KF ) holds.
In particular, degKx=(X)2 + (XKF ). For an arbitrary
divisor D on X, one can consider the subset of the field
k(X) consisting of zero and of the functions f for
which (f) + D ;;.0. This is a linear space over k of finite
dimension I(D). The dimension of the complete linear
system defined by the divisor D is I(D)-I and the cal-
culation of I(D) is an important task of the theory of
algebraic curves. The strongest relevant result is the
Riemann - Roch theorem. For smooth projective curves
this theorem is the equality
/(D)-/(K-D) = deg(D)-g+l,
where g is the genus of the curve X. If I(K - DO (or
if I(K-D)=O), one says that D is special (or, respec-
tively, non-special). For non-special divisors the
Riemann - Roch theorem yields I(D)=deg(D)- g+ l.
Each divisor of degree higher than 2g - 2 is non-special.
The class of divisors which are linearly equivalent to
a divisor D on a smooth projective curve X defines a
point on the Jacobi variety J (X) of X. This variety is
identical with the Albanese variety and with the Picard
variety of X. Points which correspond to classes of spe-
cial divisors are singular points of the Poincare divisor
on J (X). If G~ denotes the subset of points of J (X)
which correspond to the classes of divisors D with
degD =n and I(D)=r, then G~ forms a subscheme in
J(X) and
dimG~ ;;:. r(n-r+l)-(r-l)g
(the Riemann - Brill- N oether theorem). This theorem
has numerous applications, one of which will now be
described. Any divisor D for which I (D)?d defines a
rational mapping CPD of the curve X into the projective
 ALGEBRAIC CURVE

space pl(D)-I. The mapping <PD depends on the class of ponent of the unit of which coincides with the group of
D. If deg(D)~2g+ 1, then <PD defines an isomorphic points of X(k). 3) If X is a curve of genus g>l, then
imbedding of X into pm, while <PD(X) is not contained Aut(X) is always a finite group. Its order is bounded
in any proper subspace of the space pm (m I(D)-I). by the number 84(g - 1) [6]. Weierstrass points (d.
Mappings <P which correspond to a multiple nK of the Weierstrass point) on X play an important role in the
canonical class of X are the most interesting from the study of the group Aut(X) in the latter case.
point of view of the birational classification of curves. Another way of studying Aut(X) is based on the fact
If g> 1, the class 3K defines an isomorphic imbedding that all smooth projective curves are finite (ramified)
of the smooth projective curve into pSg -6. Two curves coverings of the projective line.
X and Yare birationally equivalent if and only if their Let X be a smooth projective curve defined over the
images <P3K(X) and <P3K(Y) are obtained from each field C. The set of points of the curve X(C) has the
other by a projective transformation of pSg -6. The natural structure of a one-dimensional compact analytic
study of the mapping <PK yielded a more precise charac- manifold, which is also known as a compact Riemann
teristic of curves of genus g> 1. For these curves surface. The converse is also true, i.e. any compact
<PK: X~pg-I is an isomorphic imbedding if and only if Riemann surface is obtained from some smooth projec-
X is not a hyper-elliptic curve. If <PK is an isomorphism, tive curve. Usually one uses the same symbol X for the
the curve <PK(X) is called canonical; it is defined smooth projective curve and its corresponding one-
uniquely up to projective transformations in pg -I. A dimensional complex manifold. Any connected co~plex
very important task of the theory of algebraic curves is manifol~ X can be represented as a quotient X / G,
their classification up to a birational isomorphism. A where X is a simply-connected complex manifold, an?
number of important results have been obtained in this G is a group of automorphisms of X which acts on X
field, but an adequate solution of the problem is not discretely and freely. It is noteworthy that there are
available now (1977). only three one-dimensional simply-connected connected
Smooth projective curves are subdivided into four analytic manifolds, up to an isomorphism. These are
classes: the projective line CP I (the Riemann sphere), the affine
1) curves of genus 0, birationally equivalent to PI; straight line C (the finite plane) and the interior of the
2) curves of genus I (elliptic curves), birationally unit disc D = {z: I z I < I} (the Lobachevskii plane).
equivalent to a smooth cubic curve in p2; All smooth projective curves are subdivided into three
3) hyper-elliptic curves; classes, depending on which one of the three types their
4) non-hyper-elliptic curves of 8enus g> I, biratio- universal covering belongs to.
nally equivalent to a canonical curve in pg -I (algebraic The problem of classifying smooth projective curves
curves of basic type). of a given type can be reduced to the study of discrete
The genus of a curve does not fully characterize the groups of transformations of the universal coverings
birational class of an algebraic curve. The only excep- which act freely with a relatively compact fundamental
tion are curves of genus zero. If k is the field of com- domain. In the case of the projective line, G is the iden-
plex numbers C, the set of classes of mutually- tity group; in the case of the affine straight line, G is
isomorphic elliptic curves is described by points in the isomorphic to a subgroup Q of the additive group C
quotient space H / G, where H is the upper half-plane which is a two-dimensional lattice in C; in the case of
and G is the modular group consisting of rational-linear the interior of the unit disc, G is a subgroup of motions
transformations with integral coefficients and with in the Lobachevskii plane which can be defined by
determinant equal to + 1. The space H / G has the some non-Euclidean bounded polygon. Thus, the first
structure of an analytic manifold isomorphic to C (d. class above contains a unique curve pI, the second
Elliptic curve). Classes of birationally equivalent curves class consists of complex tori C / Q that all have the
of genus g> 1 are described by points belonging to structure of a one-dimensional Abelian variety (elliptic
some algebraic variety vi( g of dimension 3g - 3, which curve), the addition of points on the torus defining the
is called the moduli variety of curves of genus g. This group structure on the respective curve. All smooth
variety is irreducible. According to one conjecture vi( g elliptic curves are obtained in this way. The field of
is unirational; however, this has been proved (by F. rational functions C(X) on an elliptic curve X~C/Q
Severi) for g< 11 only. is isomorphic to the field of meromorphic doubly-
The following results are valid for the group Aut(X) periodic (elliptic) functions with period group Q. If
of automorphisms of a smooth projective curve X. I) If f(x,y)=O is the equation of the affine model of the
X is pL then Aut(X) is the group of rational-linear curve X, then there exists a parametrization x = <j>(z),
transformations PGL(I, k). 2) If X is an elliptic curve, y = I/;{z) by elliptic functions of it (a uniforrnization of
then Aut(X) is an algebraic group, the connected com- X). The third class consists of all smooth projective

85
ALGEBRAIC CURVE
curves X of genus g> 1. In this case the field C(X) is
isomorphic to the field of meromorphic functions on D
that are invariant with respect to the group G. Such
functions are known as automorphic. Every algebraic
curve of genus g> I is uniformized by automorphic
functions (cf. Uniformization). The problem of the
classification of elliptic curves also led to a study of the
quotient D / G, but this situation is substantially dif-
ferent from the one just discussed. First, the group G
has fixed points in D; secondly, the manifold D / G is
non-compact, though containing a finite LobachevskiI
plane. The study of the general case of such groups,
and of the respective quotients, plays an important role
in modem arithmetical research.
If an algebraic curve X is defined over a non-closed
field k, one of the most important problems is that of
the existence and location of the set of rational points
X(k) of X. In the case of a smooth projective curve
over a finite field k a proof has been given of the ine-
quality I N - q -1 I :S;;;2g vq,
where N is the number of
points on X that are rational over a finite extension L
of k, q is the number of elements of Land g is the
genus of X. This inequality is equivalent to the
Riemann hypothesis concerning the zeros of the t-
function of X, viz. that all zeros of the t-function lie on
the vertical line (J = 1 / 2 (cf. Zeta-function in algebraic
geometry).
N ow let X be an algebraic curve defined over the
field of rational numbers Q. For curves of genus zero
the points of X(Q) are relatively easily found, for ellip-
tic curves the rational points constitute a finitely-
generated group (if X(Q) is non-empty), while for
curves of genus g;;. 2 there is the Mordell conjecture to
the effect that X(Q) is finite.
If the ground field k is the field of rational functions
ko(B) of a smooth projective curve B, each smooth pro-
jective curve X over k is isomorphic to the general fibre
X1J of the morphism f: V ~B of a smooth projective
algebraic surface V over ko. This morphism is uniquely
defined if it is assumed that its fibres do not contain
exclusively curves of genus 1. The set of rational points
X(k) is in one-to-one correspondence with the set of
sections V(B) of f, and X(k) is finite for curves of
genus g > 2. Curves of genus 0 and lover the field
ko(B) are studied in the theory of algebraic surfaces
(d. Elliptic surface; Ruled surface).
References
[I] SHAFAREVICH, I.R.: Basic algebraiC geometry, Springer, 1977
(translated from the Russian).
[2] WALKER, RJ.: Algebraic curves, Springer, 1978.
[3] MUMFORD, D.: Lectures on curves on an algebraic surface,
Princeton Univ. Press, 1966.
[4] CHEVALLEY, c.: Introduction to the theory of algebraiC functions
of one variable, Amer. Math. Soc., 1951.
[5] SERRE, J.-P.: Groupes algebrique et corps des classes, Hermann,
1959.
86
[6] CHEBOTAREV, N.G.: The theory of algebraiC functions,
Moscow-Leningrad, 1948 (in Russian).
[7] SPRINGER, G.: Introduction to Riemann surfaces, Addison-
Wesley, 1957.
[8] DoLGACHEV, I.V. and ISKovsKIKH, V.A.: 'Geometry of alge-
braic varieties', J. Soviet Math. 5, no. 6 (1976), 803-864. (Itogi
Nauk. i Tekhn. Algebra Topol. Geom. 12 (1974),77-170)
VE. Voskresenskil
Editorial comments. Estimate (1) above is due to G.
Castelnuovo [A1]. A proof can also be found in [A2], [A3],
[A4]. Reference [A2] also contains new results on the
Riemann-Noether-Brill theorem, e.g. it has been proven
that equality holds in the theorem for a generiC curve in the
sense of moduli; this reference also gives a survey of recent
developments in the theory of algebraic curves.
The moduli space of curves .A 9 is of general type, thus
not unirational, if 9 is odd and g~25 (J.E. Harris and D.
Mumford) or if 9 is even and g~40 (Harris), [A5].
The Mordell conjecture, i.e. every curve of genus at least
two over a number field has only a finite number of rational
points, has been proved by G. Faltings [A6].
References
[A 1) CASTELNUOVO, G.: 'Studies on the geometry of algebraic
curves', Atti R. Acad. Sci. Tarin 24 (1889), 196-223 (in
Italian).
[A2) ARilELLO, E., CORNALBA, M., GRIFFITHS, P.A. and HARRIS,
J.E.: Geometry of algebraic curves, 1, Springer, 1985.
[A3) GRIFFITHS, P.A. and HARRIs, J.E.: Principles of algebraic
geometry, Wiley, 1978.
[A4) HARTSHORNE, R.: Algebraic geometry, Springer, 1977.
[AS] MUMFORD, D. and HARRIS, J.: 'On the Kodaira dimension of
the moduli space of curves', Invent. Math. 67 (1982), 23-88.
[A6) FALTINGS, G.: 'Endlichkeitssatze fUr abelsche Varietaten Liber
Zahlk6rpern', Invent. Math. 73 (1983), 349-366.
AMS 1980 Subject Classification: 14HXX
ALGEBRAIC CYCLE on an algebraic variety - An ele-
ment of the free Abelian group the set of free genera-
tors of which is constituted by all closed irreducible
subvarieties of the given algebraic variety. The sub-
group of the group C(X) of algebraic cycles on a
variety X generated by a subvariety of codimension p is
denoted by cP(X). The group C(X) can be represented
as the direct sum
ceX) == EBcP(X).
p
The subgroup C l (X) is identical with the group of Weil
divisors (d. Divisor) on X.
In what follows X will denote a non-singular projec-
tive algebraic variety of dimension n over an algebrai-
cally closed field k. If k is the field of complex numbers
C, then each algebraic cycle Z E cP (X) defines a
(2n - 2p )-dimensional homology class
[ZlEH 2n -2p(X, Z) and, in accordance with Poincare
duality, a cohomology class y(Z)EH 2p(X, Z). The
homology (or, respectively, cohomology) classes of type
[Zl (or y(Z) are called algebraic homology (respectively,
cohomology) classes. Each analytic cycle is homologous

with an algebraic cycle. It is believed (Hodge's conjec-
ture) that an integral (2n - 2p )-dimensional cycle r on
X is homologous with an algebraic cycle if and only if
the integrals of all closed differential forms of type
(2p -q, q), q=/=p, over r are equal to zero. This conjec-
ture has only been proved for p = 1 (for n =2 [6], and
for all n [7]), for p =n -1, and for isolated classes of
varieties [4].
If W = ~ni W; is an algebraic cycle on the product of
two varieties X X T, then the set of cycles on X of the
form
is known as a family of algebraic cycles on X
parametrized by the base T. The usual requirement in
this connection is that the projection of each subvariety
W; on T be a flat morphism. If W = Wi is defined by
an irreducible subvariety, the corresponding family of
algebraic cycles on X is called a family of algebraic sub-
varieties. In particular, for any flat morphism f: X ~ Y
of algebraic varieties its fibres Xy form a family of alge-
braic subvarieties of X parametrized by the base Y. A
second particular case of this concept is that of a linear
system. All members of a family of algebraic sub-
varieties (or, respectively, algebraic cycles) of a projec-
tive variety X, parametrized by a connected base, have
the same Hilbert polynomial (respectively, virtual arith-
metic genus).
Two algebraic cycles Z and Z' on a variety X are
algebraically equivalent (which is denoted by Z "-'algZ')
if they belong to the same family, parametrized by a
connected base. Intuitively, equivalence of algebraic
cycles means that Z may be algebraically deformed into
Z'. If this definition includes the condition that the
base T is a rational variety, the algebraic cycles Z and
Z' are called rationally equivalent (which is denoted by
Z"-'ratZ'). If Z, z' EC1(X), the concept of rational
equivalence reduces to the concept of linear equivalence
of divisors. The subgroup of algebraic cycles rationally
(or, respectively, algebraically) equivalent to zero, is
denoted by C rat (X) (respectively, C alg (X. Each of
these groups is a direct sum of its components
qat(X) = n CP(X),
Crat(X)
CPrug(X) = Calg(X) n CP(X).
The quotient group Cl(X)/C~g(X) is finitely gen-
erated and is called as the N eron - Severi group of the
variety X. The problem of the quotient group
cP(X) / C~g (X) being finitely generated for p> 1
remains open at the time of writing (1977). The quo-
tient group C~g (X) / C:at (X) has the structure of an
Abelian variety (cf. Picard scheme). The operation of
intersection of cycles makes it possible to define a mul-
tiplication in the quotient group C(X) / Crat(X), con-
verting it into a commutative ring, called the Chow ring
ALGEBRAIC CYCLE
of the variety X (cf. Intersection theory).
For any Wei! cohomology theory H(X) there exists
a uniquely defined homomorphism of groups
y: cP(X) ~ H"'(X).
Two algebraic cycles Z and Z' are called homologically
equivalent (which is denoted by Z "-'homZ,) if
y(Z)=y(Z'). The subgroup of algebraic cycles that are
homologically equivalent with zero is denoted by
Chom(X). The imbedding Calg(X)CChom(X) is valid.
The quotient group C(X) / Chom(X) is finitely gen-
erated, and is a subring in the ring H (X), which is
denoted by A (X) and is known as the ring of algebraic
Wei! cohomology classes. It is not known (1986) whether
or not A (X) depends on the Weil cohomology theory
that has been chosen.
Two algebraic cycles Z and Z' are called T-equivalent
(which is denoted by Z "-'TZ') if there exists an m ~ I
such that mZ "-'algmZ'. The subgroup of algebraic
cycles that are T-equivalent to zero, is denoted by
CT(X). Two algebraic cycles Z and Z' from cP(X) are
called numerically equivalent (which is denoted by
Z"-'numZ') if the equality WZ= WZ' is valid for any
WE C n -P(X), provided both sides of the equality are
defined. The subgroup of algebraic cycles numerically
equivalent with zero is denoted by C num(X). The
imbeddings
CT(X) C Chom(X) C Cnum(X)
are valid. For divisors the groups CT(X) C1(X), n
Chom(X) n C1(X) and Cnum(X) n
C1(X) are identical
[6]. However, in accordance with the counterexample in
[5] for the case k =c
C(X) =1= Chom(X),
where Chom(X) is considered with respect to the ordi-
nary cohomology theory with rational coefficients. A
similar counterexample was established for a field k of
arbitrary characteristic and for the l-adic theory of Wei!
cohomology. The question as to the equality of the
groups Chom(X) and Cnum(X) has been solved [9].
Let X be imbedded in a projective space and let Lx
be the cohomology class of a hyperplane section. An
algebraic cohomology class
x E AP(X) = A' (X) n H"'(X)
is called primitive if xL'X -P =0. In such a case, if k is
the field of complex numbers C, the bilinear form
(a, b) ~ (-1)" Lx-"'ab
is positive definite on the subspace of primitive classes
in AP(X). A similar proposition for arbitrary k, which
is closely connected with the Wei! conjectures on the
zeta-function of an algebraic variety, has been proved
for n.;;;;2 only.
87
ALGEBRAIC CYCLE
If a variety X is defined over a field ~ that is not
algebraically closed, the Galois group G(k / k) of the
separable algebraic closure of the field k acts on the
Weil cohomology H*(X), where X=XQ!)"kk. Each ele-
ment of A' (X) is invariant with respect to some sub-
group of finite index of the group G(f / k). It is
believed (Tate's conjecture on algebraic cycles) that the
converse proposition is also true if k is finitely gen-
erated over its prime subfield. Many conjectures on the
zeta-function of algebraic varieties are based on this
assumption [2].
References
[1] BALDASSARRl, M.: Algebraic varieties, Springer, 1956.
[2] TATE, J.: 'Algebraic cohomology classes', in Summer school of
algebraic geometry Woods Hole, 1964.
[3] DoLGACHEV, LV. and ISKovsKlKH, V.A.: 'Geometry of alge-
braic varieties', J. Soviet Math. 5, no. 6 (1976), 803-864. (ftogi
Nauk. i Tekhn. Algebra Topol. Geom. 12 (1974), 77-170)
[4] KLEIMAN, S.L.: 'Algebraic cycles and the Wei! conjecture', in
Dix exposes sur la cohomologie des schemas, North-Holland,
1968, pp. 359-386.
[5] GRIFFITHS, P.A.: 'On the periods of certain rational integrals
II', Ann. of Math. (2) 90, no. 3 (1969),496-541.
[6] LEFSCHETZ, S.: L'anaiysis situs et la geometrie algebrique, Paris,
1924.
[7] HODGE, W.V.D.: The theory and applications of harmonic
integrals, Cambridge Univ. Press, 1952.
[8] 'Groupes de monodromie en geometrie,algebrique', in M. Ray-
naud, D.S. Rim and A. Grothendieck (eds.): Sem. Geom. Alg.,
Vol. 7, Springer, 1972-1973.
[9] DELIGNE, P.: 'La conjecture de Wei! 1', Pub!. Math. fHES 43
(1974), 273-308.
/. V. Dolgachev
Editorial comments. In 1983 H. Clemens proved that
CP(X) / cglg (X) is not finitely generated [A 1]. He also
proved that Chom(X) / C,(X) is not finitely generated, even
after tensoring with the field of rational numbers [A1].
A state-of-the-art survey concerning the Hodge conjec-
ture is in [A2]. See also [A3].
Much of the recent progress of the theory of algebraic
cycles is related to algebraic K-theory, see [A4].
References
[A1] CLEMENS, H.: 'Homological equivalence modulo algebraic
equivalence is not finitely generated', Publ. Math. IHES 58
(1983), 19-38.
[A2] SHIADO, T.: What is known about the Hodge conjecture,
North-Holland - Kinokuniya, 1983.
[A3] ATIYAH, M.F. and HIRZEBRUCH, F.: 'Analytic cycles on com-
plex manifolds', Topology 1 (1961),25-45.
[A4] BLOCH, S.: Lectures on algebraic cycles, IV, Duke Univ.,
1980.
AMS 1980 Subject Classification: 14C25
ALGEBRAIC DIMENSION of a topological space -
One of the different numerical invariants of the ring of
continuous functions on the given space, related to the
concept of dimension.
P.s. Aleksandrov
AMS 1980 Subject Classification: 54F45, 55M10
88
ALGEBRAIC EQUATION - An equation of the type
f" =0, where f" is a polynomial of degree n in one or
more variables (n ~o). An algebraic equation in one
variable is an equation of the form
(1)
Here n is a non-negative integer, ao, ... ,an are given,
the so-called coefficients of the equation, while x is an
unknown which has to be found. It is assumed that the
coefficients of the algebraic equation (1) are not all
equal to zero. If ao*O, then n is called the degree of
the equation.
The values of the unknown x which satisfy equation
(1), i.e. the values which, if substituted for x, will con-
vert this equation into an identity, are known as the
roots of the equation (1), or as the roots of the polynomial
j,,(x) = aOxn+alx n- 1+ ... +an- (2)
The roots of a polynomial are related to its coefficients
by Viete's formulas (cf. Viete theorem). To solve an
equation means to find all its roots contained in the
range of values of the unknown(s) under consideration.
As far as applications are concerned, the most
important case is that of coefficients and roots of an
equation that are numbers of a certain kind (e.g.
rational, real or complex). The case of the coefficients
and roots being elements of an arbitrary field may also
be considered.
If a given number (or element of a field) c is a root
of the polynomial f,,(x) then, in accordance with the
Bezout theorem, f,,(x) is divisible by x -c without
remainder. The division may be performed according to
the Horner scheme.
A number (or element of a field) c is called a root of
multiplicity k of a polynomial f (x) (where k is a non-
i,
negative integer) if f (x) is divisible by (x - c but is
not divisible by (x - ci + I. Roots of multiplicity one
are called simple roots of the polynomial, other roots are
called multiple roots.
Each polynomial f (x) of degree n >0 with coeffi-
cients in a field P has at most n roots in this field P,
each root being counted the number of times equal to
its multiplicity (consequently, there are not more than n
different roots).
In an algebraically closed field any polynomial of
degree n has exactly n roots (counted according to their
multiplicity). In particular, this statement also applies
to the field of complex numbers.
Equation (1) of degree n with coefficients from a
field P is called irreducible over P if the polynomial (2)
is irreducible over this field, i.e. cannot be represented
as the product of other polynomials of degrees lower
than n over P. Otherwise, both the polynomial and the
corresponding equation are called reducible. Polynomi-

als of degree zero and zero itself are not considered to
be reducible or irreducible. Whether a given polynomial
is reducible or irreducible over a field P depends on the
field in question. Thus, the polynomial x 2 - 2 is irredu-
cible over the field of rational numbers, since it has no
rational roots, but is reducible over the field of real
numbers: x 2 -2=(x + V2)(x - V2). Similarly, the
polynomial x 2 + I is irreducible over the field of real
numbers, but is reducible over the field of complex
numbers. Only polynomials of the first degree are
irreducible over the field of complex numbers, and any
polynomial can be decomposed into linear factors.
Only polynomials of the first degree and polynomials
of the second degree without real roots are irreducible
over the field of real numbers (and all polynomials can
be decomposed into products of linear and irreducible
quadratic polynomials). Irreducible polynomials of all
degrees exist over the field of rational numbers; exam-
ples are the polynomials of the form xn +2. The irredu-
cibility of a polynomial over the field of rational
numbers can often be established by Eisenstein's cri-
terion: If, for a polynomial (2) of degree n >0 with
integral coefficients, there exists a prime number p such
that the leading coefficient ao is not divisible by p, all
the remaining coefficients are divisible by p, and the
constant term an is not divisible by p2, then this poly-
nomial is irreducible over the field of rational numbers.
Let P be an arbitrary field. For each polynomial
f (x) of degree n> 1 that is irreducible over P, there
exists an extension of P containing at least one root of
f(x) (cf. Extension of a field); moreover, there exists a
splitting field of f (x), i.e. a minimal extension of P in
which this polynomial can be decomposed into linear
factors. Every field has an algebraically closed exten-
SIOn.
Solvability by radicals. Any algebraic equation of
degree not exceeding four is solvable by radicals. Solu-
tions of problems involving special forms of equations
of the second and third degrees were already known to
ancient Babylonians (2000 B.C.) (cf. Quadratic equa-
tion; Cubic equation). The theory of the solution of
quadratic equations was first expounded in the book
Arithmetica by Diophantus (third century A.c.). In the
sixteenth century, Italian mathematicians obtained the
solution by radicals of equations of the third and
fourth degrees with coefficients that are letters (cf. Car-
dano formula; Ferrari method). During the 300 years
which followed, fruitless efforts were made to find solu-
tions by radicals of equations of the fifth and higher
degrees with coefficients that are letters. It was finally
demonstrated by N.H. Abel in 1826 that such solutions
do not exist.
Abel's theorem in modern formulation: Let (1) be an
ALGEBRAIC EQUATION
equation of degree n >4 with coefficients a 0, . . . ,an;
let K be an arbitrary field and let P be the field of
rational functions in ao, ... ,an with coefficients from
K; then the roots of equation (1) (lying in some exten-
sion of P) cannot be expressed in terms of the coeffi-
cients of this equation by a finite number of operations
of addition, subtraction, multiplication and division
(operations which are meaningful in P) and root extrac-
tion (which is meaningful in the extension of P). In
other words, the general equation of degree n >4 is
unsolvable by radicals ([3]).
However, Abel's theorem does not contradict the fact
that some algebraic equations with numerical coeffi-
cients (or with coefficients from a given field) are solv-
able by radicals. Some special equations of degree n
are solvable by radicals (e.g. a two-term equation). A
complete solution of the problem when an algebraic
equation is solvable by radicals was given by E. Galois
about the year 1830.
The fundamental theorem on the solvability of alge-
braic equations by radicals in Galois theory can be
stated as follows. Let f (x) be a polynomial with coeffi-
cients from a field K that is irreducible over K. Then,
1) if at least one root of the equation f (x) = 0 can be
expressed by radicals in terms of the coefficients of this
equation and if the exponents of the radicals are not
divisible by the characteristic of K, then the Galois
group of this equation is solvable over K; 2) conversely,
if the Galois group of the equation f (x) = 0 is solvable
over the field and the characteristic of K is zero or is
higher than all orders of the constituent factors of this
group, then all the roots of the equation can be
represented by radicals in terms of its coefficients, all
exponents of the roots a I / n involved can be taken to
be prime numbers, and the binomial equations
xn -a =0 which correspond to these roots can be taken
to be irreducible over the fields to which these radicals
are adjoined.
This theorem was proved by Galois for the case in
which K is the field of rational numbers; in this case all
conditions involving the characteristic of the field K in
the above formulation become superfluous.
Abel's theorem is a consequence of Galois' theorem,
since the Galois group of equations of order n with
coefficients (that are letters) over the field P of rational
functions in the coefficients of the equation with coeffi-
cients from an arbitrary field K, is the symmetric group
Sn, which is unsolvable for n >4. For any n >4 there
exist equations of degree n with rational (and integral)
coefficients that are unsolvable by radicals. An example
of such an equation for n = 5 is the equation
x 5- p2 X - P = 0, where p is a prime number. In Galois
theory an algebraic equation is solved by reducing it to
89
ALGEBRAIC EQUATION
a chain of simpler equations, which are called resol-
vents (cf. Resolvent) of the original equation.
The solvability of equations by radicals is closely
connected with problems involving geometrical con-
structions with ruler and compasses; in particular, with
the problem of the division of the circle into n equal
parfs (cf. Cyclotomic polynomial; Primitive root).
Algebraic equations in one unknown with numerical
coefficients. Methods of approximate calculations (e.g.
the parabola method) are generally used to find the
roots of algebraic equations of degree higher than two
with coefficients from the field of real or complex
numbers. It is convenient to begin by getting rid of the
multiple roots. A number c is a root of multiplicity k
of a polynomial f (x) if and only if the polynomial and
its derivatives up to the order k -1, inclusive, vanish if
x =c, and if j<k)(c)=FO. If f(x) is divided by the
greatest common divisor d(x) of this polynOInial and
its derivative, then the resulting polynomial has the
same roots as f (x), but only of multiplicity one. It is
further possible to construct the polynomials whose
simple roots are all the roots of the polynomial f (x) of
equal (given) multiplicity. A polynomial has multiple
roots if and only if its discriminant vanishes.
The determination of the number of roots and
bounds on their size are frequently occurring problems.
The number
max I a; I
1+ ;>0
I ao I
can be taken as an upper bound for the modulus of
each root (both real and complex) of the algebraic
equation (1) with arbitrary complex coefficients. The
Newton method usually yields a more exact bound if
the coefficients are real. The determination of a lower
bound for the positive roots, and of upper and lower
bounds for the negative roots, are reduced to the deter-
mination of an upper bound for the positive roots.
The simplest method to determine the number of real
roots is to use the Descartes theorem. If it is known
that all the roots of a given polynomial are real (for
example, for the characteristic polynomial of a real
symmetric matrix), then Descartes' theorem yields the
exact number of roots. By considering the polynomial
f( - x), the number of negative roots of f (x) can be
found using the same theorem. The exact number of
real roots located in a given interval (in particular, the
number of all real roots) of a polynomial with real
coefficients without multiple roots can be found by the
Sturm theorem. Descartes' theorem is a special case of
the Budan - Fourier theorem, which yields an estimate
from above for the number of real roots of a polyno-
mial with real coefficients lying in a certain fixed inter-
val.
90
I t is sometimes desirable to find roots of a special
type. For instance, Hurwitz' criterion is a necessary
and sufficient condition for all the roots of an equation
(with complex coefficients) to have negative real parts
(cf. Routh- Hurwitz criterion).
There exists a method for calculating all rational
roots of a polynomial with rational coefficients. A poly-
nomial f(x) with rational coefficients has the same
roots as the polynomial g(x) with integral coefficients
obtained from f (x) by multiplication by a common
multiple of all denominators of the coefficients of f (x).
The only rational roots of a polynomial
g(x)=boxn +b)x n -) + ... +bm bn=FO, with integral
coefficients are found among the irreducible fractions
p / q in which P is a divisor of the number bn and q is
a divisor of the number b o (and only those fractions
among them such that, for any integer m, the number
gem) is divisible by p-mq). If b o=l, then all rational
roots of g(x) (if any) are integers, divisors of the con-
stant term, and can be found by trial and error.
Systems of algebraic equations. Concerning systems of
algebraic equations of the first degree, see Linear equa-
tion.
A system of any two algebraic equations of any
degree in two unknowns x and y may be written as:
f(x,y) = ao(x)yn +al(X)yn-1 + ... +an(x) = a,}
g(x,y) = bo(x)ys+bl(x)y,-l + ... +bs(x) = 0, (3)
where ai(x), b/x) are polynomials in one unknown x.
If a certain numerical value is assigned to x, a system
of two equations in one unknown y with constant coef-
ficients ai, bj is obtained. The resultant of this system is
the following determinant:
ao al an 0 0 0
0 ao an
0 0
0 0 0 ao an
R(f, g) = bo bl bs 0 0 0
0 bo bs
0 0
0 0 0 bo bs
(There are s rows of a's and n rows of b's in this deter-
minant.) The following statement is true: A number Xo
is a root of the resultant R(f, g) if and only if the poly-
nomials f(xo,Y) and g(xo,y) have a common root yo
or if the two leading coefficients ao(xo) and bo(xo) are
equal to zero.
Thus, in order to solve the system (3) one must find
all roots of the resultant R(f, g), substitute each one of
these roots into the system (3) and find the common
roots of these two equations in one unknown y. One
must also find the common roots of the two polynomi-
als ao(x) and bo(x), substitute them into (3), and verify

if the resulting equations in one unknown y have com-
mon roots. In other words, the solution of a system of
two algebraic equations in two unknowns is reduced to
the solution of one equation in one unknown and the
calculation of the common roots of two equations in
one unknown (the common roots of two or more poly-
nomials in one unknown are the roots of their largest
common divisor).
Systems of any number of algebraic equations with
any number of unknowns are solved in a similar
manner. This task involves complicated calculations. It
is connected with the so-called elimination theory.
References
[I] KUROSH, A.G.: Higher algebra, Mir, 1972 (in Russian).
[2] SUSHKEVICH, AX.: Foundations of higher algebra, Moscow-
Leningrad, 1941 (in Russian).
[3] WAERDEN, B.L. VAN DER: Algebra, 1-2, Springer, 1967-1971
(translated from the German).
[4] MANIN, YU.I.: Encyclopaedia of elementary mathematics, Vol. 4,
Moscow, 1963, pp. 205-227 (in Russian).
[5] DoMORYAD, A.P.: Encyclopaedia of elementary mathematics,
Vol. 2, Moscow, 1951, pp. 313-511 (in Russian).
/. V Proskuryakov
Editorial comments. Abel's theorem mentioned above
(the general equation of degree n >4 is unsolvable by radi-
cals) is also called the Abel- Ruffini theorem.
Good up-to- date textbooks are [A 1] and [A2].
References
[A 1] N.: Lectures in abstract algebra. Field theory and
JACOBSON,
Galois theory, 3, Freeman, 1975.
[A2] JACOBSON, N.: Basic algebra, 1-2, Freeman, 1974-1980.
AMS 1980 Subject Classification: 26CXX, 12AXX
ALGEBRAIC FUNCflON A function
y =f(x l, . . . , xn) of the variables x j, Xn that satis-
... ,
fies an equation
(1)
where F is an irreducible polynomial in y, Xl, . . . , Xn
with coefficients in some field K, known as the field of
constants. The algebraic function is said to be defined Algebraic functions of one variable. Over the field C of
over this field, and is called an algebraic function over complex numbers, an algebraic function of one variable
the field K. The polynomial F(y, Xl, . . . ,xn ) is often y = f(x) (or y(x) for short) is a k-valued analytic func-
written in powers of the variable y, so that equation (1) tion. If D (x) is the discriminant of the polynomial
assumes the form F(x,y) = Pk(X)yk+ ... +PI(x)Y+Po(x),
Pk(XIo'" ,Xn)yk+Pk_I(XIo'" ,Xn)yk-I+ ... +
+PO(XI, ... ,xn) = 0,
where Pk(Xl, ... ,xn), ... ,PO(Xl, ... ,xn) are polyno-
mials in Xl, . . . ,Xn, and with Pk(Xl,'" ,xn);;eO. The
number k is the degree of F with respect to y, and is
called the degree of the algebraic function. If k = 1, an
algebraic function may be represented as a quotient
PO(XI, ... ,xn)
y=-
PI(XI, ... ,xn)
of polynomials, and is called a rational function of
ALGEBRAIC FUNCTION
For k = 2, 3, 4, an algebraic function can
Xl, . . . ,Xn .
be expressed as square and cube roots of rational func-
tions in the variables x j, . . . , Xn; if k >4, this is
impossible in general.
The theory of algebraic functions was studied in the
past from three different points of view: the function-
theoretical point of view taken, in particular, by N.H.
Abel, K. Weierstrass and B. Riemann; the arithmetic-
algebraic point of view taken by R. Dedekind, H.
Weber and K. Hensel; and the algebraic-geometrical
point of view, which originated with the studies of A.
Clebsch, M. Noether and others (cf. Algebraic
geometry). The first direction of the theory of algebraic
functions of a single variable is connected with the
study of algebraic functions over the field of complex
numbers, in which they are regarded as meromorphic
functions on Riemann surfaces and complex manifolds;
the most important methods applied are the geometri-
cal and topological methods of the theory of analytic
functions. The arithmetic-algebraic approach involves
the study of algebraic functions over arbitrary fields.
The methods employed are purely algebraic. The theory
of valuations and extensions of fields are especially
important. In the algebraic-geometrical approach alge-
braic functions are considered to be rational functions
on an algebraic variety, and are studied by methods of
algebraic geometry (cf. Rational function). These three
points of view originally differed not only in their
methods and their ways of reasoning, but also in their
terminology. This differentiation has by now become
largely arbitrary, since function-theoretical studies
involve the extensive use of algebraic methods, while
many results obtained at first using function-theoretical
and topological methods can be successfully applied to
more general fields using algebraic analogues of these
methods.
(2)
Pk(x) '" 0,
(i.e. of the polynomial for which F(x, f(x = 0), which
is obtained by eliminating y from the equations
F( x,y ) = , aF(x, y) =
ay
to yield the equation
Pk(x)D(x) = 0,
then the roots Xl, . . . ,Xm of this last equation are
known as the critical values of y = f (x). The comple-
mentary set G=C\ {Xl, . . . ,xm } is known as the
91
ALGEBRAIC FUNCTION
non-critical set. For any point Xo EG equation (2) has k
different roots yA, ... ,y~ and the condition
of(xo,Y~) .
oy =1= 0, J=l, ... ,k,
is satisfied. According to the implicit function theorem,
in a neighbourhood of the point xo there exist k
single-valued analytic functions f6 (x), ... ,f~ (x) which
satisfy the conditions
fd(xo) = y~, F(x,fd(x =0
and which can be decomposed into a convergent series
fd(x) = y~ +a1(x - xO)+a1(x - xo)2 + . . ..
Thus, for each point Xo EG one can construct k ele-
ments of an analytic function, known as the function
elements with centre at the point Xo. For any two points
X"X2EG, any elements fl (x) andf1(x) with centres at
x, and X2, respectively, are derived from each other by
analytic continuation along some curve in G; in partic-
ular, any two elements with the same centres are also
connected in this way. If Xo is a critical point of an
algebraic function, then two cases are possible: 1) Xo is
a root of the discriminant, i.e. D(xo)=O, but
Pk(xo)=F0; or 2) Pk(xo)=O.
Case 1. Let Ko be a small disc with centre at Xo
which does not contain other critical points, and let
/J. (x), ... ,hex) be a system of regular elements with
centre at x' EK o, x' =FXO. These functions remain
bounded as x ~xo. Furthermore, let D be the circle
with centre Xo passing through x'; it is completely con-
tained inside Ko. The analytic continuation of some
given element, e.g. /J. (x), along D (in, say, the clock-
wise direction), yields an element / (x) which also
belongs to the system of elements with centre x'. This
system consists of k elements, and a minimum required
a,
finite number ~k of such turns yields the initial ele-
ment /J. (x). One obtains a subsystem /J. (x), ... ,la, (x)
of elements with centre x'; each one of these elements
may be obtained by analytic continuation of the other
by a number of turns around the point Xo; such a
subsystem is known as a cycle. Any system
flex), ... ,jk(x) can be decomposed into a number of
non-intersecting cycles
{/I (x), ... ,la, (x)},
{fa,+I(X), ... ,Ia,Ta,(X)}, ... ,
{fa,t ... +a,,+I(X), ... ,Ia,+ ... +a,(X)},
a, + ... +as=k. If a,>I, then the element /J.(x) is
not a single-valued function of x in the disc Ko, but is
a single-valued analytic function of the parameter
T=(X-XO)'/a, in a neighbourhood of T=O. In a cer-
tain neighbourhood of this point the elements
/J. (x), ... ,la, (x) of the first cycle can be represented
as convergent series
92
00 00 .
/I (x) = ~a}"; = ~a}(x-xo)'la" (4)
;=0 ;=0
00 00
Ia,(x) = ~af'''; = ~af'(x_xo)'la,;
;=0 ;=0
and similar expansions also take place for the elements
of other cycles. Such expansions of elements by frac-
tional degrees of the difference x - xo, where Xo is a
critical point, are known as Puiseux series. The
transformation T~T" r = e 2'ITj / a, which corresponds to
one tum around x 0, converts the Puiseux series of ele-
(3)
ments in one cycle into each other in cyclic order, i.e.
there is a cyclic permutation of the series and of the
corresponding elements. To turns around the critical
point correspond permutations of the elements with
centre at this point; these permutations consist of
cycles of the orders a" ... ,as, ~aj = k. The permuta-
tions defined in this way constitute the monodromy
group of the algebraic function. If at least one element
aj is greater than 1, the critical point Xo is called an
algebraic branch point of the algebraic function; the
numbers aj (sometimes aj -I) are called the branch
indices (or branch orders) of the algebraic function.
Case 2. If Pk(x)y is substituted for y, one returns to
case 1; one obtains expansions similar to (4), which
may contain a finite number of terms with negative
indices:
00 00
f(x) = ~ air ~ a;(x - xoY la. (5)
;=-P I=-P
If P >0, the point Xo is a pole of order p of the alge-
braic function. An algebraic function is usually con-
sidered on the Riemann sphere S, i.e. on the complex
plane completed by the point at infinity x = 00. The
introduction of the variable T= I / x reduces this case
to the previous case; in a neighbourhood of the point
T= (x = 00) one has the expansion
y(x) = ~ ajTj la = ~ ajX- Jla . (6)
j=-r j=-r
If r>O, then the point x = 00 is called a pole of order r.
The parameter of the expansion in the series (3), (4),
(5), (6) is called the local uniJormizing parameter for the
algebraic function. If Xo is a non-critical point of the
algebraic function, then T = x - x 0 can be taken as
parameter; if, on the other hand, Xo is a critical point,
the root (x - XO)I / a (where a is a positive integer) can
be taken as such a parameter. The population of all
elements of an algebraic function described above
forms the complete algebraic function in the sense of
Weierstrass. Algebraic functions have no singularities
other than algebraic branch points and poles. The con-
verse proposition is also true: A function y = f (x)
which is analytic, is not more than s-valued at all

points of the Riemann sphere except for a finite
number of points x I, . . . , Xm and x = 00, and has at
such points only poles or algebraic branch points, is an
algebraic function of degree k~s.
The Riemann surface of a complete algebraic func-
tion is compact and is a k-sheeted covering of the
Riemann sphere, branch points can be the critical
points and the point x = 00. Algebraic functions are the
only class of functions with a compact Riemann sur-
face. The genus of the Riemann surface of an algebraic
function is important; it is called the genus of the alge-
braic function. It can be calculated by the
Riemann - Hurwitz fonnula. The genus of a rational
function is zero, and its Riemann surface is the
Riemann sphere. The Riemann surface of an elliptic
function that satisfies a third- or fourth-degree equation
is a torus; the genus of such a function is one.
The universal covering Riemann surface of an alge-
braic function is a simply-connected two-dimensional
manifold, i.e. it has a trivial fundamental group and is
conform ally equivalent either to the Riemann sphere,
the complex plane or the interior of the unit disc. In
the first case the algebraic function is a rational, in the
second case it is an elliptic, while in the third case it is
a general function.
The unifonnization problem of algebraic functions is
closely connected with the theory of Riemann surfaces
of algebraic functions. The function y = f (x) can be
uniformized if y and x are representable as single-
valued analytic functions
y = y(t), x = x(t)
of a parameter t, which identically satisfy equation (2).
The uniformization problem is locally solved by a local
uniformizing parameter; however, it is the solution 'in
the large' that is of interest. If k = 1, i.e. if y (x) is a
rational function of x, this parameter may be the vari-
able x - x 0; if k = 2, then uniformization is attained
with the aid of a rational or a trigonometric function.
For instance, if y (x) satisfies the equation
y2_x 2 = 1,
one can take
t2 + 1
Y = -2--1'
t -
x -
or
y = sect, x = tgt.
If k = 3, 4 in the case of an algebraic function of genus
one, uniformization is achieved using elliptic functions.
Finally, if k>4 and the genus of the algebraic function
is higher than one, uniformization is realized using
automorphic functions (cf. Automorphic function).
Algebraic functions of several variables. If f is an alge-
braic function in the variables x I, . . . , x n , then t} ~ set
ALGEBRAIC FUNCTION
of all rational functions R(y, x], ... , x n ) forms a field
Kf , coinciding with the field of rational functions on
the algebraic hypersurface in (n + I)-dimensional space
defined by the equation F(y, x], ... , x n ) = O. If the
field of constants k is the field of complex numbers C
and if n = 1, then Kf is identical with the field of mero-
morphic functions on the Riemann surface of the alge-
braic function. The field Kf is an extension of finite
type of the field of constants k of transcendence degree
n (cf. Extension of a field). In par~icular, any n + 1 ele-
ments of this field are connected by an algebraic equa-
tion, so that each of them defines an algebraic function
of the remaining elements. Any extension K of finite
type of a field k of transcendence degree n is known as
an algebraic function field in n variables (or, sometimes,
as a function field). Each such field contains a purely
transcendence extension k (x I, . . . , x n ) of the field k
(called the field of rational functions in n variables). Any
element y EK satisfies some algebraic equation
I>(y, x I, . . . ,xn ) = 0, and can be considered as an alge-
braic function in the variables x I, . . . , x n . Each field K
of algebraic functions in n variables is isomorphic to
the field of rational functions on some algebraic variety
of dimension n, which is called a model of K. If the
field of constants k is algebraically closed and of
characteristic zero, then each algebraic function field
has a non-singular projective model (cf. Resolution of
singularities). Let S be the set of all non-trivial valua-
tions (cf. Valuation) of an algebraic function field K
which are non-negative on the field of constants. If
provided with the natural topology, it is known as the
abstract Riemann surface of the field K [1]. In the case
of algebraic functions in one variable, the Riemann sur-
face coincides with the set of non-singular projective
models, which in this case is uniquely defined up to an
isomorphism. Many concepts and results in algebraic
geometry on the model of a field K can be restated in
the language of the theory of valuations of fields [1],
[6]. A particularly close analogy holds for algebraic
functions in one variable, the theory of which is practi-
cally identical with the theory of algebraic curves.
Each algebraic function field in one variable is the
field of fractions of a Dedekind ring, so that many
results and concepts of the theory of divisibility in alge-
braic number fields can be applied to function fields
[12]. Many problems and constructions in algebraic
number theory motivate similar problems and construc-
tions in fields of algebraic functions and vice versa. For
instance, the application of a Puiseux expansion to the
theory of algebraic numbers led to the genesis of the p-
adic method in number theory, due to Hensel. Class
field theory, which had originally belonged to the
domain of algebraic numbers, was subsequently applied
to functions [2]. An especially close analogy exists
93
ALGEBRAIC FUNCTION
between algebraic number fields and algebraic function
fields over a finite field of constants. For instance, the
concept of a zeta-function is defined for the latter and
the analogue of the Riemann hypothesis has been
demonstrated for algebraic function fields (cf. Zeta-
fun~on in algebraic geometry).
References
[I) ZARISKl, O. and SAMUEL, P.: Commutative algebra, I, Springer,
1975.
[2) SERRE, J.-P.: Groupes algehrique et corps des classes, Hermann,
1959.
[3) SToiLOV, S.: The theory offunctions of a complex variable, 1-2,
Moscow, 1962 (in Russian and Rumanian).
[4) CHEBOTAREV, N.G.: The theory of algebraic functions,
Moscow-Leningrad, 1948 (in Russian).
[5) SHAFAREVICH, I.R.: Basic algebraic geometry, Springer, 1977
(translated from the Russian).
[6) CHEVALLEY, c.: Introduction to the theory of algebraic functions
of one variable, Amer. Math. Soc., 1951.
[7) ApPELL, P. and GOURSAT, E.: Theorie des fonctions algebriques,
1-2, Chelsea, reprint, 1976-1978.
[8) DEDEKlND, R. and WEBER, H.F.: 'Theorie der algebraischen
Funktionen einer Veranderlichen', J. Reine Angew. Math. 92
(1882), 181-290.
[9) HENSEL, K. and LANDSBERG, G.: Theorie der algebraischen
Funktionen einer Variabeln und ihre Anwendung auf algebraische
Kurven und Abelsche Integrale, Teubner, 1902.
[\0) PiCARD, E. and SIMART, G.: Theorie des fonctions algebriques
de deux variables independantes, 1-2, Chelsea, reprint, 1971.
[II) JUNG, H.W.E.: Einfohrung in die Theorir! der algebraischen
Funktionen zweier Veriinderlicher, Akademie-Verlag, 1951.
[12) HASSE, H.: Number theory, Springer, 1980 (translated from the
German).
[13) LANG, S.: A Igebraic functions, New York, 1965.
A.B. Zhizhchenko
AMS 1980 Subject Classification: 26A09, 30AXX
ALGEBRAIC GEOMETRY The branch of
mathematics dealing with geometric objects connected
with commutative rings: algebraic varieties (cf. Alge-
braic variety) and their various generalizations
(schemes, algebraic spaces, etc., cf. Scheme; Algebraic
space).
Algebraic geometry may be 'naively' defined as the
study of solutions of algebraic equations. The geometri-
cal intuition appears when every 'set of solutions' is
identified with a 'set of points in a coordinate space'.
If the coordinates are real numbers, and if the space is
two- or three-dimensional, the situation may be clearly
visualized; however, the language of geometry is also
used in more general situations. This language implies
problems, constructions and considerations which do
not follow from the point of view of pure algebra. In
turn, algebra provides a flexible and powerful
apparatus which is just as suitable for the conversion of
a tentative reasoning into a proof as for the formula-
tion of such proofs in their most obvious and most gen-
eral form.
In the case of algebraic geometry over the field of
complex numbers, every algebraic variety is simultane-
94
ously a complex-analytic, differentiable and topological
space in the ordinary Hausdorff topology. This fact
makes it possible to introduce a large number of classi-
cal structures which induce such invariants of algebraic
varieties that can be obtained by purely algebraic
means only with great difficulty or perhaps not all. The
concepts and the results of algebraic geometry are
extensively used in number theory (Diophantine equa-
tions and the evaluation of trigonometric sums), in dif-
ferential topology (both with respect to singularities
and differentiable structures), in group theory (alge-
braic groups and simple finite groups connected with
Lie groups), in the theory of differential equations (K-
theory and the index of elliptic operators), in the theory
of complex spaces, in the theory of categories (topoi,
Abelian categories), and in functional analysis
(representation theory). Conversely, the ideas and
methods of these disciplines are utilized in algebraic
geometry.
The genesis of algebraic geometry dates back to the
seventeenth century, with the introduction of the con-
cept of coordinates into geometry. Nevertheless, its
crystallization into an independent branch of science
only began in the mid-nineteenth century. The utiliza-
tion of coordinates in projective geometry created a
situation in which algebraic methods could compete
with synthetic methods. However, the achievements of
this branch of algebraic geometry still are in the
domain proper of projective geometry, but as a result it
left to algebraic geometry the traditional study of pro-
jective algebraic varieties. Modern algebraic geometry
arose as the theory of algebraic curves (cf. Algebraic
curve). Historically, the first stage of development of
the theory of algebraic curves consisted in the clarifica-
tion of the fundamental concepts and ideas of this
theory, using elliptic curves as examples. The system of
concepts and results which one now calls the theory of
elliptic curves arose as a part of analysis (rather than
geometry) - the theory of integrals of rational func-
tions on an elliptic curve. At first, these were the
integrals that received the name 'elliptic', and it was
only later that the meaning was extended to include
functions and curves (cf. Elliptic integral).
At the end of the seventeenth century, Jakob and
Johann Bernoulli noted a new interesting property of
elliptic integrals. They studied integrals expressing the
arc lengths of certain curves. They found methods for
transforming one curve into another of the same arc
length, even though the respective arcs could not be
brought into correspondence. From the analytic point
of view, this was tantamount to the conversion of an
integral into another integral, and in certain cases the
result was a transformation of an integral into itself.
e.G. Fagnano, in the first half of the eighteenth cen-

tury, gave numerous examples of such transformations.
L. Euler studied an arbitrary polynomial f (x) of the
fourth degree and posed the problem of the relations
between x and y that satisfy the equation
dx _ ~
(I)
Vj(x) - Vj(y)'
He considered (1) as a differential equation, connecting
x and y. The desired relation is the ordinary integral of
this equation. The reason for the existence of the
integral of (1) and all its special cases, discovered by
Fagnano and the Bernoulli's, is the presence of a group
law on the elliptic curve s2 = f (t), and the invariance of
the everywhere-regular differential form s -1 dt with
respect to shifts by elements of the group. Euler's rela-
tions, which interconnect x with y by (1), may be writ-
ten as
(x, Vj(xEB(e, Vj(e = (y, Vj(y,
where EB denotes the addition of points on the elliptic
curve.
Thus it is seen that these results contain both the
group law on an elliptic curve and the existence of an
invariant differential form on this curve (cf. Elliptic
curve).
After the work of Euler, the theory of elliptic
integrals was mainly developed by A. Legendre. His
considerations, beginning in 1786, are collected in a
three-volume treatise Traite des fonctions elliptiques et
integrales Euleriennes ('A treatise on elliptic functions
and Euler integrals').
The studies carried out by N.H. Abel on the theory
of elliptic functions appeared in 1827 - 1829. His start-
ing point was the elliptic integral
f
I-
0= dx
o V(l-e 2x 2)(I-e 2x 2)'
where c and e are complex numbers. Abel considered
this integral as a function of the upper limit and intro-
duced the inverse function A(8) and the function
!:l(O) = V(l-e 2A2)(I-e 2A2).
Both these functions have two periods, 2w and 2w, in
the complex domain:
I/e
w=2f dx
o V(l-c 2x2)(I-e 2x2) '
I/e
w=2f dx
o V(l-c 2x2)(l-e 2x2) '
so that the mapping x =A(8), y =t:.(8) defines a uru-
formization of the elliptic curve
y2 =(1-c 2 x 2 )(l-e 2 x 2 ) by elliptic functions.
Somewhat later than Abel, but independently, e.G.J.
Jacobi studied the function inverse to the elliptic
integral; he showed that it has two independent
ALGEBRAIC GEOMETRY
periods, and obtained several results in the problem of
transformations. By transforming Abel's expressions of
elliptic functions in series into the form of products,
Jacobi arrived at the concept of 8-functions (cf. Theta-
function) and discovered their numerous applications
not only in the theory of elliptic functions, but also in
number theory and in mechanics.
While studying the transformations of elliptic func-
tions, Abel was the first to investigate the group of
homomorphisms of a one-dimensional Abelian variety.
Finally, as a result of the publication of the inheri-
tance of e.F. Gauss and, in particular, of his diary, it
became clear that to some extent he possessed some of
these ideas long before the work of Abel and Jacobi.
The transition to the study of arbitrary algebraic curves
still took place within an analytic framework: Abel had
shown the way in which the basic properties of elliptic
integrals can be generalized to include integrals of arbi-
trary algebraic functions. Such integrals were subse-
quently named Abelian integrals (cf. Abelian integral).
The work published by Abel in 1826 became the
starting point of the general theory of algebraic curves.
It contains the concepts of the genus of an algebraic
curve and of the equivalence of divisors, and presents
equivalence criteria in terms of integrals. It leads to the
theory of the Jacobi variety of an algebraic curve.
In his dissertation, published in 1851, B. Riemann
adopted a novel principle of studying functions of a
complex variable. He assumed that such a function is
given not on the plane of the complex variable, but on
a 'many-sheeted' surface above this plane.
The surfaces introduced by Riemann (cf. Riemann
surface) are close to the modern concept of a one-
dimensional analytic manifold: A manifold on which
analytic functions are defined. Riemann posed and
solved the question of the connection between his con-
cept and the concept of an algebraic curve; the
corresponding result is nowadays called Riemann's
existence theorem. After having studied the possible
position of branching points of the surface he proved
that the set of classes depends on 3p - 3 independent
parameters if p > 1, which he called the moduli (cf.
Moduli problem).
The work of Riemann was the starting point for stu-
dies on the topology of algebraic curves; this study
explains the topological meaning of the dimension p of
the space QT[X] as one-half of the dimension of the
one-dimensional homology group of the space X (C).
Analytic methods led to the inequality
I(Dy;;:'deg(D)-p + 1. The Riemann - Roch equality
was proved by E. Roch, a student of Riemann (cf.
Riemann - Roch theorem). Finally, this study consti-
tutes the first appearance of the field k (X) as the pri-
mary object connected with a curve X, and of the con-
95
ALGEBRAIC GEOMETRY
cept of a birational isomorphism. The problem of
inversion of integrals of arbitrary algebraic functions
had been posed much earlier by Abel. Another part of
Riemann's study on Abelian functions deals with the
relations between O-functions and the inversion prob-
lem in general, in particular with the series (in p vari-
ables)
O(V) = ~eF(m)+2(m,v),
m
where m=(m), ... ,mp) runs through all integral p-
dimensional vectors,
v = (VI> ' .. ,vp ), (rn, v) = ~rnjvj,
F(rn) = ~CXjkrnjrnb CXjk = CXkj'
This series converges for all values of v if the real part
of the quadratic form F is negative definite. The princi-
pal properties of the function f) are the relations
O(v+7Tir) = OCr), O(v+CXj) = eL/VlO(v),
where r is an integer vector, aj is a column of the
matrix (ajd and L/ v) is a linear function.
It was shown by Riemann that it is possible to select
cuts a I, . . . ,ap ' b), ... ,bp which convert the surface
he has introduced into a simply-connected surface, and
to select U I, . . . ,up of the everywhere-finite integrals
on this surface such that the integrals Uj with respect to
ak are if J=I=k, and are wi if J =k, while the integrals
Uj with respect to bk form a symmetric matrix (ajd
which satisfies conditions that ensure the convergence
of the series (2), He considered the function f)
corresponding to these coefficients ajk. The periods of
an arbitrary 2n-periodic function of n variables satisfy
relations analogous to those necessary for the conver-
gence of the series which determine the f)-functions.
These relations between the periods were explicitly
written out by G. Frobenius, who showed that they are
necessary and sufficient for the existence of non-trivial
functions which satisfy the functional equation (3) (cf.
Theta-function; Abelian function). These relations are
necessary and sufficient for the existence of a mero-
morphic function with given periods which cannot be
reduced to a function in a smaller number of variables
by linear transformations. This theorem was formulated
by K. Weierstrass and proved by H. Poincare. It was
shown in 1921 by S. Lefschetz that, if the Frobenius
relations are satisfied, the f)-function defines an imbed-
ding of the variety cn / ~ into a projective space (~ is
the lattice corresponding to the given matrix of periods;
cf. Complex torus).
The concepts and the results which now form the
base of the theory of algebraic curves were created
under the influence and within the framework of the
theory of algebraic functions and their Integrals. The
purely geometrical theory of algebraic curves developed
96
as an independent branch. Thus, J. Plucker in 1834
arrived at formulas interconnecting the class of the
curve, its order and the number of its double points.
He also showed that a plane curve of order three has
nine points of inflection. However, such studies were
only of secondary importance at that time.
It was only after the studies of Riemann that the
(2)
geometry of algebraic curves became an important field
in mathematics, together with the theory of Abelian
integrals and Abelian functions. This change was
mainly due to the work of A. Clebsch. While Riemann
based his work on functions, Clebsch based his work
on algebraic curves. Clebsch and P. Gordan [10]
deduced a formula for the number p of linearly
independent integrals of the first kind (i.e. of the same
kind as the genus of the curve X), which expresses this
number in terms of the order of the curve and the
(3) number of singular points. They also showed that if
p = 0, the curve has a rational parametrization and, if
p = 1, it becomes a plane curve of order three.
An error made by Riemann proved useful in the
development of the algebraic-geometric aspect of the
theory of algebraic curves. In proving his existence
theorems, Riemann considered it self-evident that a
variational problem - the 'Dirichlet principle' - is
solvable. It was soon shown by Weierstrass that this is
not true in all cases. Accordingly, Riemann's results
remained unfounded for a certain period of time. One
way in which this difficulty could be overcome was to
prove these theorems by an algebraic method; their for-
mulation was essentially algebraic. These studies, which
were undertaken by Clebsch, greatly facilitated the
understanding of the algebraic-geometric nature of the
results of Abel and Riemann, which had hitherto been
masked by the analytic approach.
The studies initiated by Clebsch were continued and
expanded by a student of his school - M. Noether.
Noether's ideas were presented very clearly in a study
published by himself and A, Brill. They posed the
problem of the development of geometry on an alge-
braic curve lying in a projective plane as the set of
results which are invariant with respect to one-to-one
(i,e. birational) transformations (cf. Birational
geometry).
By the eighteen-fifties many special properties of
algebraic varieties of dimension higher than one
(mostly surfaces) had been discovered. For instance, a
detailed study was made of surfaces of the third degree;
in particular, it was shown by G. Salmon and A. Cay-
ley in 1849 that any cubic surface without singular
points contains 27 different straight lines. However,
these results were not connected with any general prin-
ciples for a long time, and remained unrelated to the
 ALGEBRAIC GEOMETRY

fundamental ideas in the theory of algebraic curves transformations z~(az+b)/(ez+d), where a, b, e and
which were developed at the same time. d are integers, ad - be = 1, and
The development of algebraic geometry was strongly
affected by the Italian school, in particular by L. [:~) = [~ ~) (mod 7).
Cremona, C. Segre and E. Bertini. The principal
representatives of this school were G. Castelnuovo, F. He showed that these functions uniformize the curve
Enriques and F. Severi. One of the principal achieve- x6x 1+ XTx2 + x~xo =0 of genus three. It is possible in
ments of the Italian school was the classification of this manner to distort the fundamental polygon of this
algebraic surfaces (cf. Algebraic surface). The first result group and to obtain new groups which uniformize
in this direction was reported by Bertini in 1877; he curves of genus three. A similar reasoning underlies the
gave the classification of involutive plane transforma- studies of Klein and Poincare, the latter constructing
tions or, in modern terms, a classification up to conju- automorphic functions with the aid of series which are
gation in the group of birational automorphisms of the now named after him. They both correctly conjectured
plane of all elements of order two in this group. This that any algebraic curve can be uniformized by a
classification is very simple and, in particular, it can be corresponding group, and made considerable advances
readily deduced that the quotient of the plane by a in their attempts to prove this result. The complete
group of order two is a rational surface. In other proof was only obtained in 1907 by Poincare and,
words, if a surface X is unirational and a morphism independently, by P. Koebe. A major contribution to
1: p2~X is of degree two, then X is rational. Castel- this proof was made by the studies of Poincare on the
nuovo (1893) solved (positively) the general case of the concepts of the fundamental group and a universal cov-
Liiroth problem for algebraic surfaces. He also posed ering.
and solved the problem of the characterization of The topology of algebraic curves is very simple, and
rational surfaces by numerical invariants. A classifica- was exhaustively investigated by Riemann. E. Picard
tion of surfaces was achieved by Enriques in a series of studied the topology of algebraic surfaces by a method
studies, which ended only in the first decade of the based on the study of the fibres of a morphism
twentieth century. 1: X ~p I. He studied the variation of the topology of
The principal tool of the Italian school was the study the fibre 1- I (a) with the variation of the point a Epl
of families on a surface; the families were taken linear and, in particular, conditions under which this fibre
or algebraic (the latter also being known as 'continu- contains a singular point. For instance, it was proved in
ous'). This led to the concept of linear and algebraic this way that smooth surfaces in p3 [11] are simply
equivalence. The connection between these concepts connected. Poincare also made an important contribu-
was first studied by Castelnuovo. An important contri- tion to the topology of algebraic surfaces.
bution to the subsequent development of the problem In 1921 S. Lefschetz started to use a new science -
was made by Severi. A large portion of the concepts topology - in the study of algebraic varieties over the
was defined in an analytic form, and the algebraic sig- field of complex numbers. His first aim was to obtain a
nificance was clarified in the course of time. Neverthe- simplification of the results of Poincare and Picard and
less, there are still many concepts and results which are to obtain generalizations of these results to higher-
essentially analytic, at least from the modern point of dimensional cases. However, he did considerably more
VIew. and his works opened a new domain of research in
The studies of F. Klein and H. Poincare on the prob- algebraic geometry. A further application of Lefschetz
lem of unifonnization of algebraic curves by auto- to algebraic geometry is connected with the theory of
morphic functions appeared in the early eighteen- algebraic cycles on algebraic varieties. He proved that a
eighties. Their objective was the uniformization of all two-dimensional cycle on an algebraic variety is homo-
curves by functions now known as automorphic, analo- logous to a cycle representable by an algebraic curve if
gous to the uniformization of curves of order one by and only if the regular double integral
elliptic functions. Klein's starting point was the theory jjR(x,y, z)dxdy has a zero period over this cycle.
of modular functions. The field of modular functions is Lefschetz' studies laid the foundations of the modern
isomorphic to the field of rational functions, but it is theory of complex manifolds. Such manifolds were sub-
possible to consider functions that are invariant with sequently studied by powerful tools, including the
respect to various subgroups of the modular group and theory of harmonic integrals r.yv. Hodge, G. de Rham)
to obtain more complicated fields in this manner. In and the theory of sheaves (H. Cartan, J. Leray). It
particular, Klein considered functions that are auto- could be shown by using these techniques that non-
morphic with respect to the group consisting of all singular algebraic varieties form part of an important

97
ALGEBRAIC GEOMETRY
class of complex manifolds - Kahler manifolds (cf.
Kibler manifold).
The theory of sheaves, and the related theory of vec-
tor bundles on complex manifolds, supplied a new
interpretation and considerable generalization of many
classical invariants of algebraic surfaces (of both the
arithmetic and geometric genus, and of the canonical
system). One of the most important achievements of
this theory was the creation of the theory of Chern
classes (cf. Chern class) and a considerable
generalization of the classical Riemann - Roch theorem
by F. Hirzebruch [7].
In the middle of the nineteen-twenties much work
was done on broadening the scope of algebraic
geometry from the aspect of set theory and from the
axiomatic aspect. The scope of application of algebraic
geometry was greatly extended to include complex
manifolds and algebraic varieties over arbitrary fields.
Interest in algebraic geometry over 'non-classical' fields
originated in the theory of congruences, interpreted as
equations over a finite field. It was claimed by Poincare
in his address at the International Congress of
Mathematicians in 1908 that methods of the theory of
algebraic curves may be used for studying equations
with two unknowns. The ground for a systematic con-
struction of algebraic geometry had been prepared in
the first decade of the twentieth century by the general
development of the theory of fields and the theory of
rings.
In the nineteen-thirties H. Hasse and his school
attempted to prove the Riemann hypothesis (cf.
Riemann hypotheses), which may be formulated for any
algebraic curve over a finite field; this involved the
development of a theory of algebraic curves over an
arbitrary field. The hypothesis itself was proved by
Hasse for elliptic curves. Advances in the construction
of algebraic geometry over arbitary fields are also due
to the studies of B.L. van der Waerden in the period
between 1931 and 1939. In particular, he developed
the theory of intersections in a smooth projective
variety.
A. Weil, in 1940, succeeded in formulating a proof of
the Riemann hypothesis for an arbitrary algebraic
curve over a finite field. He found two ways of proving
the hypothesis: one based on the theory of correspon-
dences of the curve X (i.e. of divisors on the surface
X X X), while the other was based on the study of its
Jacobi variety. Thus, higher-dimensional varieties are
used in both cases. Accordingly, Weil's book [5] con-
tains the construction of algebraic geometry over an
arbitrary field: the theory of divisors, cycles and inter-
sections. For the first time 'abstract' (not necessarily
quasi-projective) varieties were defined by pasting
together affine pieces. O. Zariski, P. Samuel, C. Cheval-
98
ley and J.-P. Serre introduced powerful methods of
commutative and, in particular, local algebra into alge-
braic geometry in the early nineteen-fifties.
Serre gave a definition of varieties based on the con-
cept of a sheaf. He also established the theory of
coherent algebraic sheaves, modelled on the theory of
coherent analytic sheaves, which had been introduced
only a short time before (cf. Coherent algebraic sheaf;
Coherent analytic sheaf).
In the late nineteen-fifties algebraic geometry
underwent a further radical transformation, due to
A. Grothendieck's work on the foundation of the con-
cept of a scheme. By using the language of category
theory, Grothendieck succeeded in generalizing and
clarifying many important classical constructions in
algebraic geometry, the abstract definitions of which
had been until now only little geometric. He also
founded many new branches of algebraic geometry (cf.
Abstract algebraic geometry). The language of the
theory of schemes has now become a part of the rou-
tine of modern algebraic geometry, has united it natur-
ally with commutative algebra, and has made a signifi-
cant contribution to the advances made in the study of
arithmetical problems of algebraic varieties (cf. Alge-
braic varieties, arithmetic of). The theory of schemes
directly affected one of the most important recent
achievements of algebraic geometry - the solution (by
H. Hironaka) of the problem on the existence of a
non-singular birational model of an algebraic variety
defined over a field of characteristic zero (cf. Resolution
of singularities).
References
[I] SHAFAREVICH, I.R.: Basic algebraic geometry, Springer, 1977
(translated from the Russian).
[2] HODGE, W.V.D. and !'EDOE, D.: Methods of algebraic geometry,
1-3, Cambridge Univ. Press, 1947-1954.
[3] BALDASSARRI, M.: Algebraic varieties, Springer, 1956.
[4] DIEUDONNE, J.: 'The historical development of algebraic
geometry', Amer. Math. Monthly 79 (1972),827-866.
[5] WElL, A.: Foundations of algebraic geometry, Amer. Math. Soc.,
1946.
[6] GROTHENDIECK, A. and DIEUDONNE, J.: 'Elements de
geometrie algebrique', Pub!. Math. IHES 4;8;11;17;20;24;28;32.
[7] HIRZEBRUCH, F.: Topological methods in algebraic geometry,
Springer, 1978 (translated from the German).
[8] ENRIQUEs, F.: Le superficie algebraiche, Bologna, 1949.
[9] ZARISKI, 0.: AlgebraiC surfaces, Springer, 1971.
[10] CLEBSCH, A. and GORDAN, P.: Theorie der Abelschen Funk-
tionen, Teubner, 1866.
[II] PICARD, E. and SIMART, G.: Theorie des fonctions algebriques
de deux variables independantes, 1-2, Chelsea, reprint, 1971.
Based on material from the 'Historical sketch' of [I]
Editorial comments. For another good reference on
Hirzebruch's work concerning a generalization of the classi-
cal Riemann - Roch theorem, see [A 1]. A recent historical
account of algebraic geometry is [A2].
References
[A 1] FULTON, W.: Intersection theory, Springer, 1984.

[A2] DIEUDONNE, I.: History of algebraic geometry, Wadsworth,
Monterey, 1985.
[A3] HARTSHORNE, R.: Algebraic geometry, Springer, 1977.
[A4] GRIFFITHS, PH. and HARRIs, I.: Principles of algebraic
geometry, Wiley, 1978.
AMS 1980 Subject Classification: 14-XX
ALGEBRAIC GROUP - A group G provided with the
structure of an algebraic variety in which the multipli-
cation p.: G X G~G and the inversion mapping p: G~G
are regular mappings (morphisms) of algebraic varieties.
An algebraic group is said to be defined over a field k if
its underlying algebraic variety and the morphisms p.
and p are defined over k. In such a case the set of k-
rational points of the variety G is an (abstract) group
which is denoted by G(k). An algebraic group is called
connected if its algebraic variety is connected. The
dimension of an algebraic group is the dimension of its
algebraic variety. In what follows, only connected alge-
braic groups will be considered. A subgroup H of an
algebraic group G is called algebraic if it is a closed
subvariety of the algebraic variety G. For such sub-
groups the space of co sets (left or right) can be natur-
ally provided with the structure of an algebraic variety,
defined by a universal property (d. Quotient space of
an algebraic group). If the subgroup H is also normal,
then the quotient group G / H is an algebraic group
with respect to this structure, and is calle~ an algebraic
quotient group. A homomorphism <1>: G~G of algebraic
groups is called algebraic if <I> is a morphism of their
algebraic varieties; if <I> is defined over k, it is called a
k-homomorphism. A k-isomorphism of an algebraic
group is defined in a similar manner.
Examples of algebraic groups: The general linear
group GL(n, k) (the group of all invertible matrices of
order n, with coefficients in a fixed algebraically closed
field k); the group of triangular matrices; an eUiptic
curve.
There are two main types of algebraic groups, with
altogether different properties: Abelian varieties and
linear algebraic groups (d. Abelian variety; Linear alge-
braic group). The type of a particular group is deter-
mined exclusively by properties of its variety. An alge-
braic group is called an Abelian variety if its algebraic
variety is complete. An algebraic group is called linear
if it is isomorphic to an algebraic subgroup of a general
linear group. An algebraic group is linear if and only if
its algebraic variety is affine. These two classes of alge-
braic groups have a trivial intersection: If an algebraic
group is both an Abelian variety and a linear group,
then it is the identity group. The study of arbitrary
algebraic groups reduces to a great extent to the study
of Abelian varieties and linear groups. In particular, an
arbitrary algebraic group contains a unique normal
ALGEBRAIC GROUP OF TRANSFORMATIONS
linear algebraic subgroup H such that the quotient
group G / H is an Abelian variety [AI]. Numerous
examples of algebraic groups which are neither linear
algebraic groups nor Abelian varieties are given by the
theory of generalized Jacobi varieties for algebraic
curves with singular points [3] (cf. Jacobi variety). A
natural extension of the class of algebraic groups leads
to the concept of a group scheme.
References
[I] BOREL, A.: Linear algebraic groups, Benjamin, 1969.
[2] MUMFORD, D.: Abelian varieties, Oxford Uniy. Press, 1974.
[3] SERRE, I.-P.: Groupes algebrique et corps des classes, Hermann,
1959.
[4] CHEYALLEY, c.: 'La theorie des groupes algebriques', in Proc.
Internat. Congress Mathematicians Edinburgh, 1958, Cambridge
Uniy. Press, 1960, pp. 53-68.
[5] DEMAZURE, M. and GROTHENDIECK, A.: 'Schemas en groupes',
in Sem. Geom Alg. 1963-1964, Leet. Notes in Math., Vol. 151-
153, Springer, 1970. B B TJ" k
. . "en ov
v.P. Platonov
Editorial comments.
References
[A1] CHEVALLEY, c.: 'Une demonstration d'un theoreme sur les
groupes algebriques', J. Math. Pures Appl. 39 (1960), 307-
317.
AMS 1980 Subject Classification: 14KXX, 20GXX
ALGEBRAIC GROUP OF TRANSFORMATIONS -
An algebraic group G acting regularly on an algebraic
variety V. More precisely, it is a triplet (G, V, T) where
T: G X V ~ V (r(g, x) = gx) is a morphism of algebraic
varieties satisfying the conditions: ex =x, g(hx)=(gh)x
for all x E V and g, hE G (where e is the unit of G). If
G, V and T are defined over a field k, then (G, V, T) is
called an algebraic group of k-transformations. For
instance, (G, G, T), where T is the adjoint action or an
action by shifts, is an algebraic group of transforma-
tions. If G is an algebraic subgroup in GL(n) and T is
its natural action on the affine space V = k", then
(G, V, T) is an algebraic group of transformations. For
each point x E V one denotes by G (x) = {gx: g E G} the
orbit of x, and by Gx = {gEG: gx =x} the stabilizer of
x. The orbit G (x) need not necessarily be closed in V,
but closed orbits exist always, e.g. orbits of minimal
dimension are closed. An algebraic group of transfor-
mations is sometimes understood to mean a group G
which is acting rationally (but not necessarily regularly)
on an algebraic variety V (this means that T: G X V ~ V
is a rational mapping, and the above properties of T are
valid for ordinary points). It was shown by A. Well [3]
that there always exists a variety VI, birationally iso-
morphic to V, and such that the action of G on VI
induced by the rational action of G on V is regular.
The problem of describing the orbits, stabilizers, fields
of invariant rational functions (cf. Invariants, theory of),
and of constructing quotient varieties are fundamental
99
ALGEBRAIC GROUP OF TRANSFORMATIONS
in the theory of algebraic groups of transformations
and have numerous applications.
References
[1) BoREL, A.: Linear algebraic groups, Benjamin, 1969.
[2A) DIEUDONNE, J.A. and CARREll, J.B.: Invariant theory, old and
new, Acad. Press, 1971.
[2B) MUMFORD, D.: Geometric invariant theory, Springer, 1965.
[2C) MUNFORD, D.: 'Projective invariants of projective structures
and applications', in Proc. Internat. Congress Mathematicians
Stockholm, 1962, Inst. Mittag-Leffler, 1%3, pp. 526-530.
[2D) SESHADRI, C.S.: 'Quotient spaces modulo reductive groupes
and applications to moduli of vector bundles on algebraic
curves', in Actes du Congres Internat. Mathematiciens Nice,
197JJ, Yol. 1, Gauthier-Yillars, 1971, pp. 479-482.
[3) WElL, A.: 'On algebraic groups and homogeneous spaces',
Amer. J. Math. 77, no. 2 (1955), 355-391.
v.P. Platonov
Editorial comments. The notion in question is also
called an algebraic transformation space.
AMS 1980 Subject Classification: 20GXX, 20HXX
ALGEBRAIC INDEPENDENCE - A concept in the
theory of field extensions (cf. Extension of a field). Let
K be some extension of a field k. The elements
b 1> ,bn EK are called algebraically independent over
k if for each polynomial f(x) , ... ,xn ) with coefficients
from k which is not identically equal to zero,
feb), ... ,bn)*O. Otherwise, the elements b), ... ,bn
are called algebraically dependent. An infinite set of ele-
ments is called algebraically independent if each one of
its finite subsets is algebraically independent; otherwise
it is called algebraically dependent. The definition of
algebraic independence may be extended to the case
where K is a ring, and k a sub ring [I].
Algebraic independence of numbers. Complex numbers
a), ... ,an are called algebraically independent if they
are algebraically independent over the field of algebraic
numbers, i.e. if for any polynomial P(x), ... ,xn ) with
algebraic coefficients, not all of which are zero, the
relationship P(a), ... , an)*O is valid. Otherwise,
a), ... ,an are called algebraically dependent. The con-
cept of algebraic independence of numbers is a general-
ization of the concept of transcendency of a number
(the case n = I). If several numbers are algebraically
independent, each one of them is transcendental. It is
usually very difficult to prove that given numbers are
algebraically independent. The eXlstmg analytic
methods of the theory of transcendental numbers give a
solution of this problem for the values of certain classes
of analytic functions. For instance, it was found that
the values of the exponential function e Z , the argu-
ments of which are algebraic and linearly independent
over the field of rational numbers, are algebraically
independent. A similar result was obtain~d for Bessel
functions (cf. Siegel method). Several general theorems
were also established concerning algebraic indepen-
dence of the values in algebraic points of E-functions
100
which satisfy linear differential equations with coeffi-
cients from the field of rational functions [2], [3]. Alge-
braic independence has been proved for numbers aP
and at, where a*O, I is an algebraic number and f3 is
a cubic irrationality. A number of theorems have also
been proved which deal with the algebraic non-
expressibility of numbers; a concept close to that of
algebraic independence.
It is possible to impart a quantitative facet to the
qualitative concept of algebraic independence if the
measure of algebraic independence (cf. Algebraic
independence, measure of) of these numbers is con-
sidered. The analytic methods referred to above yield
estimates from below for the measure of algebraic
independence of certain classes of numbers. General
theorems treating the estimation of the measure of
algebraic independence of the values of E-functions
have been established [3].
References
[I) LANG, S.: Algebra, Addison-Wesley, 1974.
[2) FEL'DMAN, N.!. and SHIDLOvsKII, A.B.: 'The development and
present state of the theory of transcendental numbers', Russian
Math. Surveys 22, no. 3 (1967), 1-79. (Uspekhi Mat. Nauk 22,
no. 3 (1967), 3-81)
[3) SHIDLOVSKIi, A.B.: 'On arithmetic properties of values of ana-
lytic functions', Trudy Mat. Inst. Steklov. 132 (1973), 169-202
(in Russian).
A.B. ShidlovskiF
Editorial comments. The above-mentioned theorem on
the algebraic independence of the values of e Z is known as
the Lindemann- Weierstrass theorem. Schanuel's
conjecture, 'transcendence degree
(x 1, . . . ,Xn , eX' , ... ,eX,),;:, n for any Q-linear independent
X1, .. ,XnEC', is still (1986) unproved. An analogue of the
Lindemann - Weierstrass theorem for values of the Weier-
strass Ie -function has been proved recently [A1], [A2] by
using Gel'fond - Schneider-type arguments and zero esti-
mates for polynomials on algebraic group varieties. Another
method in algebraic independence is Mahler's method for
values of functions satisfying functional equations [A3].
References
[A 1] WDSTHOLZ, G.: 'Ueber das abelsche Analogon des Lin-
demannsche Satzes I " Inv. Math. 72 (1983), 363-388.
[A2] PHILIPPON, P.: 'Varietes abeliennes et independence
algebrique II', Inv. Math. 72 (1983),389-405.
[A3] LOXTON, J. and POORTEN, A.1. VAN DER: 'Transcendence and
algebraic independence by a method of Mahler', in A Baker
and D. Masser (eds.): Transcendence theory, Acad. Press,
1977, Chap!. 15.
AMS 1980 Subject Classification: 10F37
ALGEBRAIC INDEPENDENCE, MEASURE OF -
The measure of algebraic independence of the numbers
aI, ... ,am is the function
cIJ(al,'" ,am;n,H) = min 1 P(al,. ,am) I,
where the minimum is taken over all polynomials of
degree at most n, with rational integer coefficients not

all of which are zero, and of height at most H. For
more details see Transcendency, measure of.
A.B. Shidlovskit
AMS 1980 Subject Classification: 12AXX
ALGEBRAIC IRRATIONALITY - An irrational alge-
braic number. An algebraic irrationality is the root of
an irreducible polynomial of a degree at least two, with
rational coefficients.
A.I. Galochkin
AMS 1980 Subject Classification: 12AXX
ALGEBRAIC K-THEORY - A branch of algebra,
dealing mainly with the study of the so-called K-
functors (Ko. K], etc., cf. K-functor); it is a part of
general linear algebra. It deals with the structure theory
of projective modules and their automorphism groups.
To put it more simply, it is a generalization of results
obtained on the existence and uniqueness (up to an
automorphism) of a basis of a vector space and other
group-theoretical facts concerning linear groups over
fields. On passing from a field to an arbitrary ring R
these theorems usually become invalid, and the
Grothendieck group Ko(R) and the Whitehead group
K] (R) are, in a certain sense, a measure of their devia-
tion from being true. Similar generalizations of the
structure theorems of linear algebra appear in topology.
A vector space can be regarded as a special case of a
vector bundle. These objects may be studied with the
aid of the homotopy theory of vector bundles and of
topological K-tbeory. It is important to note in this
connection that a projective module can be regarded as
the module of sections of a vector bundle. This explains
the choice of the class of projective modules as the
object of the theory. Algebraic K-theory makes exten-
sive use of the theory of rings, homological algebra,
category theory and the theory of linear groups.
Algebraic K-theory has two different historical ori-
gins, both in the field of geometry. The first is related
to certain topological obstructions. The starting point
was the introduction of the concept of Whitehead tor-
sion, which is connected with the homotopy
equivalence of finite complexes and is an element in the
Whitehead group, the latter being some quotient group
of the group K] (R), where R = ZII is the integral group
ring of the fundamental group II. The next step con-
cerned topological spaces X that are dominated by a
finite complex, and their generalized Euler characteris-
tic X(A), which is an element of group Ko(Zw). The
computation of the Whitehead group and L-groups
(which is, strictly speaking, an algebraic problem
concerning group rings), was in fact one of the first
objectives of algebraic K-theory. Both K2 and other
higher functors have topological applications of the
ALGEBRAIC K-THEORY
same type (for example, an obstruction to the deforma-
tion of a pseudo-isotopy of a closed manifold into an
isotopy lies in some quotient group of the group
K2(ZW. Algebraic studies of the Whitehead group
began in the nineteen-fourties. A related field is the
study of the structure of linear groups over arbitrary
rings, in particular, the theory of determinants over a
skew-field [10].
The second origin of algebraic K-theory was an alge-
braic proof of the Riemann - Roch theorem [7] and its
generalizations by A. Grothendieck in 1957. These con-
siderations involved the introduction of the K-functor
K(X) as the group of values of a universal additive
functor on coherent sheaves on a smooth algebraic
variety. Moreover, the previously familiar representa-
tion rings, Witt rings (cf. Witt ring) of classes of qua-
dratic forms, etc., turned out to be related construc-
tions. The K-functor was then transferred to topology,
in which it found numerous applications, and with its
help several previously unsolved problems could be
dealt with.
It became clear, moreover, that this construction
reveals new perspectives in the understanding of old
analytical problems (the index problem of elliptic
operators), topological problems (extraordinary homol-
ogy theories), and the theory of group representations.
However, the development of algebraic K-theory for
rings (beginning with the establishment of the
correspondence (analogy) between projective, finitely-
generated modules and vector bundles) was hindered
by the fact that an adequate concept, analogous to that
of suspension in topology, was lacking in algebra.
The nineteen-fifties and nineteen-sixties saw the
beginning of the systematic study of projective modules
over finite groups, and the development of one of the
most important ideas on which algebraic K-theory is
based - the idea of 'stabilization', the essence of
which, roughly speaking, is that general relationships
are more clearly manifested on passing to the limit of
the dimension of the objects studied (e.g. linear groups
or projective modules). Connections were noted
between algebraic K-theory and the reciprocity laws of
the theory of algebraic numbers and algebraic func-
tions; studies were made of problems connected with
congruence subgroups (cf. Congruence subgroup) and
an algebraic analogue of the Bott periodicity theorem
- the theory of polynomial extensions - was
obtained.
For a ring R with a unit element, the Grothendieck
group Ko(R) is defined as the Abelian group generated
by the isomorphism classes of finitely-generated projec-
tive R-modules, with the defining relation:
101
ALGEBRAIC K-THEORY

where [P] is the class of modules isomorphic to the If fz is also an epimorphism, then the sequence is sup-
module P. Let GL(n, R) be the general linear group plemented by the terms
over R, let A ~ II~ ?II K 2 (R) ~ K 2 (R 1)ffiK2 (R 2 ) ~ K 2 (R') ~ K1(R) ~ .. '.

be the imbedding of GL(n, R) in GL(n + 1, R), let
GL(R) be the direct limit of the groups GL(n, R), and
let E(R) be the subgroup in GL(R) generated by the
elementary matrices et, i.e. by the matrices that have
an element hER at the (i, j)-th place and agree with
the unit matrix in all other places. E(R) then coincides
with the commutator of GL(R). The quotient group
GL(R)/ E(R) is denoted by KI(R), and is known as
the Whitehead group. Finally, the Steinberg group
St(n, R) for n;;;.3 is defined by the generators XA.'J' hER ,
l~i,j~n, i=/=j, and the relations
xAx"" = x A+I'.
'J'
lJ I)
[xt, xjd = x>;t if i=/=I;
If I is a two-sided ideal of R, then the Mayer- Vietoris
sequence makes it possible [8] to define the relative
functors K;(R, I), which yield an exact sequence
K 2(R,J) ~ KiR) ~ K 2 (R/ J) ~ K1(R,J) ~ K1(R) ~
~ K1(R / J) ~ Ko(R, J) ~ Ko(R) ~ Ko(R / J).
A fairly complete study has been made of the
behaviour of K-functors on passing from a ring R to its
localization with respect to a central, multiplicatively-
closed system. In particular, if certain conditions on R
are satisfied, then the following exact sequence has
been found for the functor Go(R):
U Go(R / (s ~ Go(R) ~ GO(S-I R) ~ O.
seS

[xt,x~tl =1 if }=/=k, i=/=I. If R is commutative, Ko(R) becomes a ring with a unit

Passing to the direct limit, one obtains the group St(R) element by introducing the multiplication induced by
and a natural homomorphism the te~sor product of modules. There exists a split epi-
morphism of Ko(R) onto the ring H (R) of continuous
</>: St(R) ~ E(R),
integer-valued functions (the ring Z is given the discrete
with topology) on the spectrum of R (cf. Spectrum of a ring).
<f>(xt) = et
The kernel of this_homomorphism is denoted by Ko(R).
The kernel kercf> is denoted by K 2(R) (the Milnor
It is known that Ko(R) is the nil radical of Ko(R) and,
group). It coincides with the centre of St(R). Thus,
if R is Noetherian and if the dimension of its maximal
~o, KI and K2 are functors from the category of rings
spectrum is a, thel! Ko(Rt+ 1 =0. If this dimension is
mto the category of Abelian groups. Each of the func-
at most 1, then Ko(R) is isomorphic to the Picard
tors K 0 and K I can be characterized as a functor from
group Pic(K).
finitely-generated projective modules to Abelian groups
For arithmetical rings there are finiteness theorems
that satisfies certain properties and is universal with
for the functors K;(R) and G;(R). In fact, if A is the
respect to these properties. Such a 'universal' character-
ring of integers or the ring of polynomials over a finite
ization makes it possible to define analogues of the
field, and R is an R -order and at the same time an R-
functors Ko and KI on 'sufficiently good' categories. In
lattice in a semi-simple finite-dimensional algebra over
particular, for the category of Noetherian R-modules
the field of fractions of a ring A, then the groups K;(R)
functors G;(R) quite close to K;(R) can be defined.
and G;(R) are finitely generated (i =0, 1).
Examples oj the groups K;(R). Let R be a skew-field
The development of algebraic K-theory was stimu-
and let R* be its multiplicative group. Ko(R)=Z is
lated by studies carried out on the problem of
then the group of integers, KI(R)=R* /[R*,R*]; and
congruence subgroups: Do all subgroups of finite index
K 2 (Z) is the cyclic group of order two. If R is a finite
in an arithmetical group contain some congruence sub-
field, then K2(R)=0.
group? This question is closely connected with the
An important result in algebraic K-theory is the
problem of computing the group K I (R, J) for ideals J
exact Mayer- Vietoris sequence for a Cartesian square.
in R.
The diagram below represents a Cartesian square of
. ~f the results concerning the stable structure of pro-
ring homomorphisms in which JI is an epimorphism;
JectIve modules one can mention the following
h,
R ~ R2
theorem: If R is a commutative Noetherian ring whose
maximal spectrum has dimension d, and A is a
hd th,
R ~R'
module-finite R-algebra, then any finitely-generated
I /,
projective A-module P such that
then there is an exact sequence Pm ~ A~1+1 ffiQ
K1(R) ~ K 1(R 1)ffiK 1(R 2 ) ~ K1(R') ~
for all maximal ideals m of R is isomorphic to A EEl N
~ Ko(R) ~ K o(R 1)ffiK o(R 2 ) ~ Ko(R'). (here M m is the localization of the module M at m).

102

Another important theorem on the structure of projec-
tive modules is the cancellation theorem: Let R, A and
the module P be as above. Let Q be a finitely-
generated projective A-module, and let M and N be
arbitrary A-modules. Then it follows from
QtBPfJJM-QfJJN that
PfBM ~ N.
The stable rank of a ring R is closely connected with
problems of the stable structure of projective modules.
Thus, if R is a commutative ring of stable rank smaller
than d, then
K1(R)::: GL(d, R)/ E(d, R).
In connection with the theory of induced representa-
tions of groups, the functors K; for group rings have
been studied. One of the results of these studies is that
if G is a finite group of order n and C is the set of
cyclic subgroups of G, then the index of the subgroup
:s Im(K;(Rc )--,>K;(RG
CEC
in K;(RG) is divisible by n if i =0, 1,2.
Regarding polynomial ring extensions it
that if R is a regular ring, then
Ko(R[t)) ~ Ko(R[t, t -1)) ::: Ko(R),
K 1(R[t)) ~ K1(R).
Moreover, the sequence
0----> K 1(R) ----> K 1(R[t))fBK 1(R[t- I ))---->
----> KI(R[t, t-I)) --'> Ko(R) ----> 0
is exact for any ring R.
One result in the computation of the functor K 2 (R)
is the theorem of Matsumoto: If R is a field, then K 2 (R)
is given by the generators lea) (which are in one-to-one
correspondence with all the non-zero elements a of R)
and the relations l(a)l(l-a)= 1 for a:;i=l.
In the nineteen-seventies there appeared numerous
versions of the definitions of the functors K; for i ;;;;.2.
It has been shown [9] that these theories coincide and
yield the classical functors Kn if n ~2. In several cases
effective methods of computation for higher K-groups
were found. The development of unitary K-theory ([9],
Vol. 3), which studies analogous problems for modules
on which quadratic and bilinear forms are defined, also
began in that decade.
References
[I] ATIYAH, M.F.: K-theory: lectures, Benjamin, 1967.
[2] BASS, H.: Lectures on topics in algebraic K-theory, Tata Inst.
Fundam. Res., 1966.
[3] BASS, H.: Algebraic K-theory, Benjamin, 1968.
[4] SWAN, R.G.: Algebraic K-theory, Springer, 1968.
[5] SWAN, R.G. and EVANS, E.G.: K-theory offinite groups and
orders, Springer, 1970.
[6] THOMAS, c.B. and Moss, R.M.F. (EDS.): Algebraic K-theory
and its geometric applications, Springer, 1969.
[7] MANIN, Yu.!.: 'Lectures on the K -functor in algebraic
ALGEBRAIC LOGARITHMIC SINGULAR POINT
geometry', Russian Math. Surveys 24, no. 5 (1969), 1-89.
(Uspekhi Mat. Nauk 24, no. 5 (1969), 3-86)
[8] MiLNOR, J.: Introduction to algebraic K-theory, Princeton Univ.
Press, 1971.
[9] BASS, H. (ED.): Algebraic K-theory (Battelle Inst. Conf), 1-3,
Springer, 1973.
[10] ARTIN, E.: Geometric algebra, Interscience, 1957.
A. V. Mikhalev
A.I. Nemytov
Editorial comments. Algebraic K-theoretic ideas and
results have become most important in certain parts of func-
tional analysis centering around C'-algebras (cf. C'-
algebra). Especially in the form of KK-theory (or Kasparov
K-theory). Cf. e.g. [A2].
In algebraic geometry there are important connections
with the Chow groups (cf. Chow ring).
References
[A1] MAGURN, B. (ED.): Reviews in K-theory 1940-1984, Amer.
Math. Soc., 1985.
[A2] CuRTZ, J.: 'K-theory and C'-algebras', in A. 8ak (ed.): Alge-
braic K-theory, number theory and analysis, Springer, 1984,
pp.55-79.
AMS 1980 Subject Classification: 18F25
IS known
ALGEBRAIC LATIICE, compactly-generated lattice -
A lattice each element of which is the union (i.e. the
least upper bound) of some set of compact elements (d.
Compact lattice element). A lattice is isomorphic to the
lattice of all sub algebras of some universal algebra if
and only if it is both complete and algebraic. These
conditions are also necessary and sufficient for the lat-
tice to be isomorphic to the congruence lattice of some
universal algebra (the Gratzer-Schmidt theorem). In
both cases it is assumed that the arity of the operations
of the universal algebra is finite.
References
[1] BIRKHOFF, G.: Lattice theory, Amer. Math. Soc., 1967.
[2] GRATZER, G.: Lattice theory, Birkhauser, 1978.
T.S. Fofanova
Editorial comments. Algebraic lattices are an important
special case of continuous lattices, for which see [A1]. The
original reference to the Gratzer-Schmidt theorem is [A2].
References
[A1] GIERZ, G., HOFMANN, K.H., KEIMEL, K., LAWSON, J.D.,
MISLOVE, M.V. and SCOrf, D.S.: A compendium ot continu-
ous lattices, Springer, 1981.
[A2] GRATZER, G. and ScHMIDT, E.T.: 'Characterizations of
congruence latlices of abstract algebras', Acta. Sci. Math.
(Szeged) 24 (1963), 34-59.
AMS 1980 Subject Classification: 06A23
ALGEBRAIC LOGARITHMIC SINGULAR POINT -
An isolated singular point z 0 of an analytic function
f (z) such that in a neighbourhood of it the function
f (z) may be represented as the sum of a finite number
of terms of the form
(z - zo)-S[ln(z - zo)t g(z),
103
ALGEBRAIC LOGARITHMIC SINGULAR POINT
where s is a complex number, k is a non-negative
integer, and g(z) is a regular analytic function at the
point z 0 with g(z oh6 0.
References
[ll BIEBERBACH, L.: Analytische Fortsetzung, Springer, 1955,
Sect. 3.
E.D. Solomentsev
AMS 1980 Subject Classification: 30B40, 30BXX
ALGEBRAIC NUMBER - A complex (sometimes,
real) number that is a root of a polynomial
f(x) = a.x' + ... +alx +ao
with rational coefficients, not all of which are zero. If a
is an algebraic number, then, among all polynomials
with rational coefficients and a as a root, there exists a
unique polynomial <p(x) of lowest degree' with leading
coefficient equal to one, which is therefore irreducible
(d. Irreducible polynomial). It is called the irreducible,
or minimal, polynomial of the algebraic number a. The
degree n of the minimal polynomial <p(x) is also called
the degree of the algebraic number a. The existence of
irreducible polynomials of any degree n implies the
existence of algebraic numbers of degree n. All rational
numbers, and only such numbers, are algebraic
numbers of the first degree. The number i is an alge-
braic number of the second degree, since it is a root of
the polynomial x 2 + 1, while 21/ n, where n is any posi-
tive integer, is an algebraic number of degree n, being a
root of the irreducible polynomial xn - 2.
The roots aI, ... ,an of the irreducible polynomial
are called the conjugates of a, and are also algebraic
numbers of degree n. All numbers conjugate with a are
distinct. Apart from its degree, another important
characteristic of an algebraic number is its height,
which is the analogue of the denominator of a rational
fraction. The height of an algebraic number a is the
greatest absolute value of the coefficients of the irredu-
cible and primitive polynomial with integral rational
coefficients that has a as a root (d. Primitive polyno-
mial). The sum, the difference, the product and the
quotient of two algebraic numbers (except for division
by zero) are algebraic numbers; this means that the set
of all algebraic numbers is a field. A root of a polyno-
mial with algebraic coefficients is an algebraic number.
An algebraic number is called an algebraic integer if
all the coefficients of its minimal polynomial are
rational integers. For instance, i and 1+ V2 are alge-
braic integers, being roots of the polynomials x 2 + 1
and x 2 -2x -1 .
The concept of an algebraic integer is a generaliza-
tion of the concept of a rational integer (a rational
integer m is an algebraic integer, being the root of the
polynomial x - m). Many properties of rational
integers are also displayed by algebraic integers. Thus,
104
the algebraic integers form a ring; on the other hand,
the real algebraic integers form an everywhere-dense set
in R, while the rational integers form a discrete set. It
was shown in 1872 by G. Cantor that the set of all
algebraic numbers is denumerable, and this implied the
existence of transcendental numbers (cf. Transcendental
number).
A root of any (not necessarily irreducible) polyno-
mial with rational integer coefficients and leading coef-
ficient one is an algebraic integer. Moreover, a root of a
polynomial with algebraic integer coefficients and with
(1) leading coefficient one is an algebraic integer. In partic-
ular the k-th degree root of an algebraic integer is an
algebraic integer. For any algebraic number a there
exists a positive integer r such that ra is an algebraic
integer (in analogy with rational numbers). The smal-
lest possible number r is the modulus of the leading
coefficient of the irreducible primitive polynomial with
rational integer coefficients that has a as a root. All
conjugates of an algebraic integer are also algebraic
integers.
One says that an algebraic integer f3 is divisible by
an algebraic integer a, a*O, if there exists an algebraic
integer y such that f3 = ya. Many divisibility characteris-
tics of rational integers also hold for algebraic integers.
An algebraic integer t: is called an algebraic unit (or
unit for short) if it is a divisor of 1, that is, if 1 / t: is an
algebraic integer. A unit is a divisor of any algebraic
integer. The inverse of a unit is a unit; the numbers
conjugate to a unit are units; all divisors of a unit are
units; a product of a finite number of units is a unit.
An algebraic integer is a unit if and only if the product
of all its conjugates is -+-1. The k-th roots of unity are
units, each one having modulus 1. There exists an
infinite set of other units which do not have modulus 1.
For instance, the numbers 2 - v'3 and 2 + v'3 are
units, being the roots of the polynomial x 2 -4x + 1.
Moreover, their powers constitute units which may be
arbitrary large or small. The field of rational numbers
contains only two units - + 1 and - 1.
Two algebraic integers are called associated if they
differ by a factor which is a unit. There is another
important difference between the ring of algebraic
integers and the ring of rational integers. The concept
of an irreducible integer (in analogy to a prime
number) cannot be introduced into the former. This
may be seen from the fact that a root of an algebraic
integer is an algebraic integer. The concept of an
irreducible number (apart from the class of associated
numbers) can be introduced in certain subfields of the
field of all algebraic numbers, the so-called algebraic
number fields. It turns out, however, that the decompo-
sition of an algebraic integer into irreducible factors is
not always unique.

Algebraic numbers cannot be very closely approxi-
mated by rational and algebraic numbers (Liouville's
theorem). It is this fact which led in 1844 to a proof of
the existence of transcendental numbers. The problem
of approximation of algebraic numbers by rational
numbers is one of the more difficult problems in
number theory; attempts to solve it yielded very impor-
tant results, including the Thue, Thue - Siegel and
Thue - Siegel- Roth theorems, but its ultimate solution
is still nowhere in sight. Another very difficult problem
is the expansion of algebraic numbers into continued
fractions. Real algebraic numbers of the second degree
(quadratic irrationalities) can be represented as infinite
periodic continued fractions. Nothing is known so far
(the nineteen-seventies) about the expansion of real
algebraic numbers of degree at least three into ordinary
continued fractions.
A complex number is called an algebraic number over
a field P C C if it is a root of a polynomial (1) with
coefficients from P. The minimal polynomial, the
degree over P and the conjugate numbers over P for
algebraic numbers over P are defined in a similar
manner. A root of a polynomial with coefficients that
are algebraic numbers over P is an algebraic number
over P.
An algebraic number of arbitrary degree n does not
necessarily exist over any field P. For instance, only
algebraic numbers of the first degree exist over the field
of complex numbers: the numbers of the field them-
selves. A given algebraic number can be of different
degrees with respect to different fields; thus, the
number i is an algebraic number of the second degree,
but it is of the first degree with respixt to the field of
complex numbers. The set of all algebraic numbers
over a field P forms a number field.
Algebraic numbers, and algebraic number fields,
were first systematically studied by e.F. Gauss (Gauss-
ian numbers of the form a + bi, where a and bare
rational numbers). Gauss developed the arithmetic of
Gaussian integers as a base for the theory of biqua-
dratic residues. In their study of the theory of cubic
residues, e.GJ. Jacobi and F.G. Eisenstein created the
arithmetic of numbers of the form a +bp, where
p=( -1 + v=3) / 2 is a cubic root of unity, while a
and b are rational numbers. His attempts to produce a
proof of the Fennat theorem led E. Kummer to con-
duct a deep study of cyclotomic fields (cf. Cyclotomic
field), to introduce the concept of an ideal, and to
create the elements of the theory of algebraic numbers.
The theory of algebraic numbers was further developed
by P. Dirichlet, L. Kronecker, D. Hilbert and others.
Russian mathematicians - E.!. Zolotarev (theory of
ideals), G.F. Vorono! (cubic irrationalities, units of
cubic fields), A.A. Markov (cubic fields), Yu.V.
ALGEBRAIC NUMBER THEORY
Sokhotskii (theory of ideals) and others - also made
significant contributions.
The concept of an algebraic number and the related
concept of an algebraic number field are very impor-
tant ideas in number theory and algebra. Algebraic
numbers, which are a generalization of rational
numbers, form subfields of algebraic numbers in the
fields of real and complex numbers with special alge-
braic properties. The development of the theory of
algebraic numbers greatly influenced the creation and
development of the general theory of rings and fields.
Algebraic numbers have found numerous applica-
tions in various branches of number theory, algebra
and other branches of mathematics: the theory of
forms, Diophantine equations, Diophantine approxima-
tions, transcendental numbers, geometry of numbers,
algebraic geometry, Galois theory, etc.
References
[1] TSCHEBOTAROW, N.G. [N.G. CHEBOTAREV): Grundzuge der
Galois'schen theorie, Noordhoff, 1950 (translated from the Rus-
sian).
[2) HECKE, E.: Lectures on the theory of algebraic numbers,
Springer, 1981 (translated from the German).
[3) LANDAU, E.: Einflihrung in die elementare und analytische
Theorie der algebraischen Zahlen und der Ideale, Chelsea,
reprint, 1949.
[4) LANDAU, E.: Vorlesungen iiber Zahlentheorie, 3, Hirzel, 1927.
[5) LANG, S.: AlgebraiC number theory, Addison-Wesley, 1970.
A.B. Shidlovski{
Editorial comments.
References
[A1] WEISS, E.: Algebraic number theory, McGraw-Hili, 1963.
[A2] WElL, A.: Basic number theory, Springer, 1967.
AMS 1980 Subject Classification: 12AXX
ALGEBRAIC NUMBER TIlEORY - The branch of
number theory with the basic aim of studying proper-
ties of algebraic integers in algebraic number fields K
of finite degree over the field Q of rational numbers
(cf. Algebraic number).
The set of algebraic integers OK of a field K / Q -
an extension K of Q of degree n (cf. Extension of a
field) - can be obtained from an integral basis
(WI, ... ,wn ); this means that each algebraic integer
(i.e. element of OK) can be written in the form
XIWI + ... +xnwn where all the Xi run through the
rational integers (i.e. Z). Moreover, such a representa-
tion is unique for each algebraic integer in K.
However, properties o{rational integers often do not
have obvious analogues for algebraic integers. The first
such property is related to units, the invertible elements
of OK (cf. Unit). The field of rational numbers has only
+ 1 and - 1 as units, but a general algebraic number
field may contain an infinite number of units. E.g. con-
sider the real quadratic field Q( Vii), where D > 1 is a
rational integer not equal to a square. If, moreover,
105
ALGEBRAIC NUMBER THEORY
D;;E 1 (mod 4), then (1, YD) is an integral basis for it.
=
The PeU equation x 2 - Dy2 I has an infinite number
of solutions2,Y)EZXZ. Any of them gives rise to a
unit x +Y YD of Q( YD). In fact,
(x+yVDXx-yVD) = I,
and
IVD = x-yVD
x+y D
is also an algebraic integer in Q(YD). The units of
Q(YD) form an infinite multiplicative group (the group
of Pell units). The question arises: What is the structure
of this group?
A second property of rational numbers without an
obvious analogue for algebraic numbers is related to
the theorem on unique factorization of rational integers
n in prime factors:
n = p~1 ... p%'.
For algebraic numbers such a unique factorization need
not hold. E.g. consider the field Q(v=5). In it, the
number 6 has two essenti, different factorizations:
6=23, 6=(l+v=5)(1- -5). For extensions of
higher degree the situation becomes still more compli-
cated. The question arises: What becomes of the
unique factorization theorem, does it have a meaning at
all in algebraic number fields?
A third property without an obvious analogue is
related to prime numbers. An ordinary prime number
need not remain a prime number in an algebraic
number field. E.g. in the field Q(-v=--l) of Gaussian
numbers the prime number 5 has a factorization:
5=(2+ -v=--l)(2- -v=--l). However, the prime
number 7 remains prime in this field. The question
arises: Do there exist general laws governing the
behaviour of prime numbers in algebraic number fields
of higher degree? In other words, is it possible to find
rules that would give a definite answer to the following
question: Does a given prime number remain prime in
a field K / Q or does it split in it, and if it splits, in
how many factors?
Finally, a last (fourth) problem concerns the general
structure of algebraic number fields. The field Q is the
minimal field of characteristic zero without proper sub-
fields. Any other algebraic number field does have sub-
fields. E.g., Q is a subfield of any algebraic number
field. The question arises: How many subfields are con-
tained in a given extension K / Q, finitely or infinitely
many, and what is their structure?
These are four main problems in algebraic number
theory, and answering them constitutes the content of
algebraic number theory. It is quite natural to start
with the fourth problem, since its answer will shed light
on the other three. The problem was solved by E.
Galois in the eighteen-twenties (d. Galois theory). The
106
fact that the number of subfields of an extension K / Q
of degree n over Q is finite follows from the existence
of a one-to-one correspondence (the fundamental Galois
correspondence) between all subfields of the normal clo-
sure of K (cf. Nonnal extension) and the subgroups of
its Galois group, a finite group of finite order (number
of elements) at most n!.
The structure of the group of units of a field was elu-
cidated by P. Dirichlet. His basic idea can be given
by taking the group of Pell units (cf. above) as exam-
ple. Any power of such a unit (both positive or nega-
tive) is a unit. There is a fundamental unit
11=Xo +yo YD, and all other units are plus or minus
integral powers of it; hence, the Pell units form the
product of < - 1> with an infinite cyclic group. This
fact is a special case of the general Dirichlet unit
theorem for algebraic number fields: If the degree of a
field K / Q is n =r1 +2r2, where r1 is the number of
real imbeddings K~R and r2 the number of complex-
conjugate pairs of complex irnbeddings K ~C, then the
group of units U of K is the product of a finite cyclic
group C and r=r1 +r2 -I infinite cyclic groups:
U = ex <'II!> X ... X <'IIr>.
Here 111, ... ,11, are independent fundamental units
and C is the group of roots of unity contained in K.
The norm of any unit in a field, i.e. the product of the
unit by its conjugates, is equal to + I.
The problem of the non-unique factorization of alge-
braic integers in algebraic number fields was solved by
E. Kummer, who started from a special case; he tried
to solve Fermat's last theorem on the impossibility of
solving the equation x P+yP =zP in non-zero integers
for any prime number p > 2. Kummer expanded the
left-hand side using p-th roots of unity, and hence the
problem was reduced to one about the algebraic
integers of Q(t), f a primitive p-th root of unity. If
there would be unique factorization into prime factors
in Q(f). then it would have been sufficient to prove
that not all prime factors at the left-hand side occur
with an exponent that is a multiple of p. This was
Kummer's first point of view, but Dirichlet pointed out
to him that unique factorization need not hold. In
order to overcome this difficulty Kummer introduced
ideal numbers, thus altering the entire structure of alge-
braic number theory for the future. The concept of an
ideal number arises from the fact that if a field k does
not contain prime numbers into which any algebraic
integer in k can be split, then there is a field K / k of
finite degree over k in which there does exist a collec-
tion of numbers that play the role of primes for k.
These numbers were called ideal by Kummer (since
they do not lie in the original field k). By introducing
ideal numbers the theorem on unique factorization in k
 ALGEBRAIC NUMBER THEORY

holds. Two numbers in a field that differ by a unit I(x) =Jil (x) .. .t:.~(x) (modp).
(so-called associated numbers) have one and the same The corresponding factorization is called Kummer's for-
ideal factors. (Note that ideal numbers are defined rela- mula (or Kummer's decomposition)
tive to the original field k; for another field k' one
(1)
must construct an extension K' / k' (of possibly dif-
ferent degree over k') in which all ideal numbers of k' where ~j are prime ideals in K=k(O).
are contained.) In principle, this equation solves the third problem in
Kummer also introduced the concept of the ideal algebraic number theory, but it is a local equation in
class number of a field k. Two ideal numbers are said to the sense that it is necessary to check it for each prime
belong to the same class if their quotient belongs to k. ideal separately. The problem of partitioning the set of
The number of classes thus obtained is called the ideal all prime ideals in classes such that the factorization
class number of k. He obtained the following important law is the same in any given class and such that, more-
result: The class number h of k is finite, and the classes over, it is possible to give a simple description of these
form an Abelian group under multiplication. Thus, any classes, is solved by class field theory for extensions
ideal number can be regarded as the h-th root of some K / k with Abelian Galois group Gal(K / k).
element of the original field k. The class number can be Equation (1) yields a preliminary concept of a class.
explicitly described in terms of other field constants Let n be the degree of the extension K / k and let f; be
(the regulator, the discriminant and the degree of the the relative degree of the ideal ~j. Computing the rela-
field). tive norms N K / k of both parts of (1) leads to

Subsequently, the concept of an ideal number was (2)

replaced by that of an ideal, equivalent to it, which can
in which aj and f; are natural numbers. For n fixed
be described in terms of the field k itself. Already in equation (2) has a finite number of solutions, so that
the nineteen-fifties the concept of an ideal was general-
the set of all prime ideals of k can be partitioned into a
ized to the more comprehensive concept of a divisor. finite number of classes and one can collect the prime
Therefore, the modern theory of Kummer may be ideals whose Kummer decomposition corresponds to
stated in terms of divisors. For algebraic number one solution (aj,f;) of (2) into one class. The number of
fields, however, the concepts of a divisor and an ideal prime ideals with the property that some aj > 2, is fin-
coincide, and such fields only will be considered in the ite, and they all are divisors of the discriminant ~ of
sequel. The concept of an ideal is related to that of K / k. Only infinite classes are of interest, so that
non-associated numbers, which helps in understanding classes for which aj;;;;.2 can be ignored.
the deep connections between Kummer's theory and In order to simplify the exposition, the field K / k is
Dirichlet's theory of units. Although Kummer did not considered to be normal from now on. In such fields
succeed in solving Fermat's problem, his ideas extended the condition
far beyond this problem and the concept of an ideal 11= ... = 1m
has now become fundamental in many branches of holds. Therefore, the set of all tJ not dividing ~ is parti-
mathematics. tioned into d(n) classes, where d(n) denotes the number
Since prime ideal numbers are defined relative to a of divisors. The class with m = n is of special interest, in
field or, in modern terminology, since they are prime it
ideals in a field, the third problem on the factorization 11==lm=l.
of ordinary prime numbers in an algebraic number A prime ideal tJ in it has maximal number of prime
field can be stated in the following general form. Sup- divisors in K / k:
pose one is given a field k and a prime ideal tJ of its P = PI ... Pn (3)
ring of algebraic integers. The question is: Does tJ
remain prime in an extension K=k(O) or does it split Such tJ are said to split completely (or to be totally split),
in a product of prime ideals of the ring of algebraic and their class is called the principal class of k relative
integers of K'? In the latter case: By which law does it to K / k. It is a principal object of study in class field
split? This question led to class field theory, a central theory. In order to be able to define the principal class
part of modern algebraic number theory. The first solu- via (3) it is necessary to give a proof of the facts that
tion of this problem was given by Kummer, who prime ideals satisfying (1) do exist in k and that there
proved that if 0 is a root of an irreducible polynomial are infinitely many such ideals. Therefore a basic prob-
f(x) and 1,0, ... ,fI' -) is an integral basis for OK over lem in class field theory is to define the principal class
Ob then tJ splits in k(O) 'by the same law as f(x), does in terms of the field k itself in such a way that its infin-
in the residue field (modtJ). In other words, the factori- ite nature would follow. This problem has been com-
zation of tJ in k(O) is determined by the congruence pletely solved for Abelian extensions K / k.

107
ALGEBRAIC NUMBER THEORY
In order to be able to understand the ideas of class
field theory better, it is necessary to introduce the gen-
eral concept of an ideal class group. Kummer's defini-
tion given above then corresponds to the modem con-
cept of an absolute ideal class group. The general
concept of an ideal class group which is used nowadays
is due to H. Weber and T. Takagi, [5].
Weber introduced the concept of the conductor of a
class group. Let m be an integral ideal of k, let Lm be
the subgroup of principal ideals (a) of k given by a=l
(mod m) and let A m be the subgroup of all ideals of k
that are relatively prime with m. A subgroup H m satis-
fying Lm ~Hm ~Am can be regarded as a group of
principal ideals, and a (generalized) class group
Am / H m is constructed in this way. For m = 1 and
HI = L lone obtains Kummer's definition. In the gen-
eral case H m consists of progressions (mod m) whose
residue classes form a subgroup of the whole multipli-
cative group (mod m). The order h m of this generalized
class group satisfies the estimates 1:S;;;;h m :S;;;;h'cp(m),
where h is the order of the absolute class group and <I>
is Euler's phi-function. For different conductors ml and
m2 the class groups may be equivalent if H m and H m,
l
consists in fact of the same progressions (mod/), where
f=l.c.m.(ml, m2)' If one agrees not to distinguish
between equivalent class groups, then one obtains the
concept of a class group according to Weber with a con-
ductor f that is the greatest common divisor of all
equivalent conductors. A class field according to Weber
is a field K / k in which only the prime ideals in its
principal class Hf , and only these, split completely.
Dirichlet's theorem on prime ideals in progressions,
which is valid in every field k, implies that there are
infinitely many prime ideals. Weber has shown that in
certain particular cases the Galois group Gal(K / k) of
a class field and the class group Af / Hf of k are iso-
morphic.
A new point of view on class field theory, which is
still valid, originated with D. Hilbert. He understood
that between all relatively Abelian extensions of a field
k and all class fields for this field there must exist a
one-to-one correspondence. This correspondence can be
stated as follows. If for some conductor f one con-
structs the Weber class group, then there is only one
normal extension K / k in which the prime ideals of
Weber's principal class, and only these, will split com-
pletely. Moreover, the Galois group of K / k is iso-
morphic to the Weber class group and the discriminant
D of K / k consists of the same prime ideals as the con-
ductor f The converse is also true: If an Abelian exten-
sion with Galois group Gal(K / k) is give.n, then there
exists a method (subsequently explicity formulated by
Takagi) by which one can uniquely construct a princi-
pal class Hf such that the class group Af / Hf is iso-
108
morphic to Gal(K / k), only the prime ideals of Hf
split completely in K, and such that the conductor f
has the same prime divisors as the discriminant D of
K / k. The 'duality' was stated by Hilbert in 1900 as a
hypothesis (he proved it only in special cases). In its
general form it was proved by Takagi.
The next important stage of development of class
group theory is related to the name of E. Artin, who
revealed the special role of the canonical isomorphism
between the Galois group and the ideal class group. He
proved that this isomorphism is given by the Frobenius
automorphism op of an Abelian extension K / k,
defined as
Here N p is the absolute norm of the ideal p, a runs
through all numbers of K and p is a prime ideal of k.
The automorphism op (nowadays denoted by
[-T]]
depends (in an appropriate group) only on the ideal
class to which p belongs. It is multiplicative:
[ \: k] [\: k] = [~l ~: ],
where the symbol at the right-hand side is understood
as an automorphism of the class to which lJIlJ2 belongs.
With this in mind, Artin introduced the symbol
[~]
on the entire group Af of ideals 0 of k that are rela-
tively prime with the conductor f It is called the Artin
symbol and realizes an isomorphism between the class
group Af / Hf and the Galois group Gal(K / k), whose
explicit from is expressed in Artin's reciprocity law:
[~] = 1
if and only if oEHf (reciprocity as a correspondence
between the groups Gal(K / k) and Af / Hf ). From
this one can obtain the classical form of the reciprocity
law in the language of Kummer's power reciprocity
symbol (one has to consider the field K=k(a 1 / n ), d.
Reciprocity laws). This form, in tum, implies the
reciprocity law for Hilbert's symbol. In all three forms
the reciprocity law is regarded as an isomorphism of
groups, and the symbols of Artin, Kummer and Hilbert
are regarded as group elements realizing this isomor-
phism. However, each of them also has a numerical
value, which is equal to some n-th root of unity. There-
fore one can formulate the reciprocity law as follows.
Given the value of (a / {3)n, it is required to find the
value of the symbol (a / f3)n reciprocal to it, i.e. to exhi-
bit the explicit form of the function of (a, {3) defined by

In this form the law first appeared in the work of CF.
Gauss (cf. Gauss reciprocity law) for quadratic fields
and in the work of Kummer for cyclotomic fields of
prime degree. The study of the reciprocity law in this
form was subsequently continued by CG.J. Jacobi,
F.G.M. Eisenstein, Hilbert, H. Hasse, and others, but
they only obtained partial results. I.R. Shafarevich [6]
solved this problem in its general form in 1948, basing
himself on the idea of establishing a connection
between the arithmetic definition of the symbol
[]
and the analytic definition of an Abelian differential
(ld/1 on a Riemann surface. He gave a construction of
[]
which precisely corresponds to the determination of the
residue of (ld/1 in the .):J-adic field. For this purpose he
introduced E-functions, also called Shafarevich func-
tions, in terms of which he also obtained an explicit
form of the reciprocity law in the general case.
In the late nineteen-twenties Hasse introduced the
idea of doing class field theory for one prime ideal of k
at a time, and reformulated and proved many theorems
for an Abelian extension Kp of a local field kp- At first
the idea was of secondary importance only (the results
were obtained as a constX!uence of the local theory),
but at the end of the nineteen-thirties C Chevalley
proved the extraordinary importance of the local point
of view in the creation of class field theory. He intro-
duced the concept of an idele group and formulated
general class field theory from the local point of view.
Following this the local-global principle became esta-
blished in class field theory. Subsequently it was
extended and refined (cf. [5]), and as a result Abelian
class field theory took a structured and completed
form. The problem of creating non-Abelian class field
theory for normal extensions with non-Abelian Galois
group remains.
The exposition above relates mainly to the qualitative
aspects of algebraic number theory. In questions of
quantitative estimation and methods algebraic number
theory is intimately connected with analytic number
theory. It is also based, to a large extent, on properties
of r-functions and L-functions of algebraic fields.
References
[I] BOREVICH, Z.I. and SHAFAREVICH, I.R.: Number theory, Acad.
Press, 1966 (translated from the Russian).
[2] POSTNIKOV, M.M.: Galois theory, Moscow, 1962 (in Russian).
[3] WEYL, H.: AlgebraiC theory of numbers, Princeton Univ. Press,
1959.
[4] LANG, S.: Algebraic number theory, Addison-Wesley, 1970.
ALGEBRAIC NUMBER THEORY
[5] CASSELS, J.W.S. and FROHLICH, A. (EDs.): Algebraic number
theory, Acad. Press, 1967.
[6] SHAFAREVICH, I.R: 'General reciprocity laws', Mat. Sb. 26
(68), no. 1(1950),113-146 (in Russian).
[7] 'Hilbert's problems', Bull. Amer. Math. Soc. 8 (1902), 437-479
(translated from the German).
[8] ARTIN, E. and TATE, J.: Class field theory, Benjamin, 1967.
[9] WElL, A.: Basic number theory, Springer, 1974.
[10] LANG, S.: AlgebraiC numbers, Addison-Wesley, 1964.
[II] BoURBAKI, N.: Commutative algebra, Addison- Wesley, 1966
(translated from the French). A.!, Vinogradov
Editorial comments. An algebraic number field of
degree n is a degree n extension of the field a of rational
numbers. An element aEK is an algebraiC integer if it is
integral over Z, i.e. if it is a root of a monic polynomial in
Z[XJ . These form a free Abelian subgroup of rank n, called
the ring of integers OK of K and also the principal order of
K; a basis of OK is an integral basis, also called a minimal
basis.
For Kummer's approach to the Fermat problem cf. [A3].
The link (equivalence) between Kummer's ideal numbers
and ideals comes from the fact that for the prime ideals in
OK there is a field extension L / K in which they become
principal ideals.
The relative degree 1; of an ideal \l3; is the degree of the
residue field extension OL / \l3; over OK / p where
n
P=\l3; OK'
For more details, definitions, newer results, and more pre-
cise statements of much of the above cf. various more
specified articles in this Encyclopaedia. In particular the arti-
cles Algebraic number; Unit; Extension of a field; Regula-
tor of an algebraic number field; Discriminant; Frobenius
automorphism (also for Artin symbo~, Kummer theorem
on decomposition of ideals, and especially class field
theory.
The ideal class groups Am / Hm are also often known as
ray class groups mod m, and the corresponding fields as ray
class fields. The class field corresponding to the ray class
group A1 / L1 is the Hilbert class field. The definitions
above ignore some problems with the Archimedean primes;
in particular the prinCipal ideals (a) making up Lm must
have a totally positive, i.e. such that all real conjugates of a
are >0.
The discussion of class field theory above (up to the men-
tion of ideles) is in terms of the older ideal-theoretic formula-
tion as in e.g. Hasse's Zahlbericht [A1], which in fact can
be advantageously consulted for more details and precise
statements. For a discussion of the relations between the
ideal-theoretic formulations and the modern idele based
class field theory [A2] is recommended.
Meferences
[A 1] HASSE, H.: Bericht Dber neuere Untersuchungen und Prob-
leme aus der Theorie der algebrai'sche Zahlkorper, 1.11, Phy-
sika Verlag, 1965.
[A2] NEUKIRCH, J.: Class field theory, Springer, 1986, Chapt. 4,
Sect. 8.
[A3] POLLARD, H.: The theory of algebraic numbers, Math. Assoc.
Amer., 1950.
[A4) WEISS, E.: Algebraic number theory, McGraw-Hili, 1963.
AMS 1980 Subject Classificatiol"l: 12-XX
109
ALGEBRAIC OPERATION
ALGEBRAIC OPERATION, n-ary operation, on a set
A - A mapping
of the n-th Cartesian power of the set A into the set A
itself. The number n is known as the arity of the alge-
braic operation. Historically, the concepts of binary
(n =2) and unary (n = 1) operations were the first to be
considered. Nullary (n =0) operations are fixed ele-
ments of the set A; they are also known as distinguished
elements or constants. In the twentieth century the con-
cept of an infinitary operation appeared, i.e. a mapping
w: A a ~A, where a is an arbitrary cardinal number. A
set with a system of algebraic operations defined on it
is called a universal algebra.
T.M. Baranovich
Editorial comments. The study of infinitary operations
actually started in the late nineteen-fifties [A 1]. A nullary
operation is also called a noughtary operation [A2].
References
[A1] STOMINSKl, J.: 'The theory of abstract algebras with infinitary
operations', Rozprawy Mat. 18 (1959).
[A2] COHN, P.M.: Universal algebra, Reidel, 1981, 13-14.
AMS 1980 Subject Classification: 08A40
ALGEBRAIC POLYNOMIAL OF BEST APPROXI-
MATION - A polynomial deviating least from a given
function. More precisely, let a measurable function
f(x) be in lp[a, b] (p;;;:'l) and let Hn be the set of alge-
braic polynomials of degree not exceeding n. The quan-
tity
(*)
is called the best approximation, while a polynomial for
which the infimum is attained is known as an algebraic
polynomial of best approximation in lp[a, b]. Polynomi-
als which deviate least from a given continuous func-
tion in the uniform metric (p = 00) were first encoun-
tered in the studies of P.L. Chebyshev (1852), who con-
tinued to study them in 1856 [1]. The existence of alge-
braic polynomials of best approximation was esta-
blished by E. Borel [2]. Chebyshev proved that P~(x) is
an algebraic polynomial of best approximation in the
uniform metric if and only if Chebyshev alternation
occurs in the difference f (x) - P~(x); in this case such
a polynomial is unique. If p> 1, the algebraic polyno-
mial of best approximation is unique due to the strict
convexity of the space Lp- If P = 1, it is not unique, but
it has been shown by D. Jackson [3] to be unique for
continuous functions. The rate of convergence of
En{j)p to zero is given by Jackson's theorems (cf.
Jackson theorem).
In a manner similar to (*) an algebraic polynomial of
best approximation is defined for functions in a large
number of unknowns, say m. If the number of variables
m;;;:' 2, an algebraic polynomial of best approximation
in the uniform metric is, in general, not unique.
110
References
[1] CHEBYSHEV, P.L.: Complete collected works, Vol. 2, Moscow,
1947, pp. 478; 152-236 (in Russian).
[2] BOREL, E.: Lefons SW' les fonctions de variables reelles et les
developpements en series de polynomes, Gauthier-Villars, 1905.
[3] JACKSON, D.: 'A general class of problems in approximation',
Amer. J. Math. 46 (1924),215-234.
[4] GARKAVI, A.L.: 'The theory of approximation in norrned linear
spaces', ftogi Nauk. Mat. Anal. 1967 (1969),75-132 (in Rus-
sian). Yu.N. Subbotin
Editorial comments. Instead of the long phrase 'alge-
braic polynomial of best approximation' one also uses the
shorter phrase 'best algebraic approximation', which is not
to be confused with the phrase 'best approximation' for the
least error En{ f)p.
References
[A 1] ACHIEZER, N.I. [N.I. AKmEZER]: Theory of approximation,
Ungar, 1956 (translated from the Russian).
[A2] CHENEY, E.W.: Introduction to approximation theory, Chel-
sea, reprint, 1982.
[A3] MEINARDus, G.: Approximation von Funktionen und ihre
numerische Behandlung, Springer, 1964.
AMS 1980 Subject Classification: 41 A 10, 41 A50
ALGEBRAIC SPACE - A generalization of the con-
cepts of a scheme and an algebraic variety. This gen-
eralization is the result of certain constructions in alge-
braic geometry: Hilbert schemes, Picard schemes,
moduli varieties, contractions, which are often not real-
izable in the category of schemes and require its exten-
sions. However, the category of algebraic spaces is
closed under these constructions, which renders the
algebraic space a natural object of algebr~c geometry.
Any scheme S defines some sheaf S in the etale
topology of the category of schemes, which in tum
uniquely defines the scheme S. An algebraic space is a
sheaf of sets F in the etale topology of schemes satisfy-
ing the condition of local representability (in the etale
topology): There exists a scheme U and a sheaf mor-
e.hism U~F such that for a~y sc~eme V and morphism
V~F the fibred product uXFV is represented by a
scheme Z, and the induced morphism of schemes Z ~ V
is a surjective etale morphism. The scheme U is known
as the etale covering of the sheaf F, which is a quotient
sheaf of t~e sh~af U according to the etale equivalence
relation UXFU. This last proposition reveals the
geometric sense of an algebraic space as a quotient
scheme according to the etale equivalence relation.
Morphisms of an algebraic space are defined as mor-
phisms of sheaves; the category of schemes becomes
identical with a complete subcategory of the category
of algebraic spaces.
Many concepts in the theory of schemes are applica-
ble to algebraic spaces: point, local ring, etale topology,
Zariski topology, field of functions, structure sheaf and
coherent sheaves. Many results in the theory of schemes

such as Serre's affinity criterion (cf. Affme scheme), and
the theorem of finiteness and existence of a proper
morphism, can be applied to algebraic spaces.
Finer results include the representability of Picard
and Hilbert functors in the category of algebraic
spaces. If a flat equivalence relation is given on an alge-
braic space, then the factorization by this relation
yields an algebraic space (such a situation occurs, for
example, when there is a free action of a finite group
on the space). Finally, an algebraic space permits a
contraction of a subspace with an ample conormal
sheaf.
All algebraic spaces contain a subspace that is open
and dense in the Zariski topology and that is a scheme.
One-dimensional and non-singular two-dimensional
algebraic spaces are schemes, but this is not true of
three-dimensional or singular two-dimensional algebraic
spaces; a group in the category of algebraic spaces over
a field is a scheme. Complete algebraic spaces of
dimension n over the field of complex numbers have
the natural structure of a compact analytic space with n
algebraically-independent meromorphic functions.
References
[1] ARTIN, M.: Algebraic spaces, Yale Univ., 1969.
[2] KNurSON, D.: AlgebraiC spaces, Springer, 1971.
VI. Danilov
Editorial comments. The notion of an algebraic space
was introduced by M. Artin and, in a slightly different but
equivalent form, it was studied by B. Moishezon [A 1]. The
theorem cited at the end of the article (on subspaces) is
due to him.
References
[A 1] MOISHEZON, B.: 'AlgebraiC varieties and compact complex
spaces', in Proc. internat. congress mathematicians Nice,
1970, Gauthiers-Villars, 1971, pp. 643-648.
AMS 1980 Subject Classification: 14JXX, 14A20
ALGEBRAIC SURFACE - A two-dimensional alge-
braic variety. Together with the class of algebraic curves
(cf. Algebraic curve), the class of algebraic surfaces is
the class of algebraic varieties which has been most
thoroughly studied. The richness of problems and the
variety of ideas used to solve them makes the theory of
algebraic surfaces one of the most interesting fields of
algebraic geometry.
Unlike algebraic curves, two birationally isomorphic
algebraic surfaces need not be biregularly isomorphic.
During the classical period of the development of alge-
braic geometry (1868 - 1920) algebraic surfaces were
understood to mean a class of birationally equivalent
surfaces, and its representations were considered to be
imbedded in complex three-dimensional projective
space Cp3, in which they were defined by a single alge-
braic equation f(xo, Xl> X2, X3)=0. The foundations of
ALGEBRAIC SURFACE
the theory of algebraic surfaces were laid by A. Clebsch
and M. Noether (ca. 1870), who defined the first
important invariants of algebraic surfaces: the
geometric genus and the canonical class. The subse-
quent development of the theory of algebraic surfaces is
due to the Italian school of algebraic geometers: C.
Segre, L. Cremona, E. Bertini, G. Castelnuovo, F.
Enriques, F. Severi, and others. Topological and tran-
scendental methods were introduced into the theory of
algebraic surfaces by E. Picard, H. Poincare and S.
Lefschetz. Many general facts arrived at in these studies
proved to be valid for algebraic varieties and schemes
of arbitrary dimension as well, while other results up
till now have not been generalized to higher-
dimensional cases. The latter include the classification
of algebraic surfaces by numerical invariants, criteria of
rationality and ruledness of algebraic surfaces, and the
theory of minimal models. Most of these results were
critically reviewed and demonstrated once again in the
nineteen-fifties and nineteen-sixties using modem coho-
mological methods (1], [5], [8].
General features of algebraic surfaces. Unless otherwise
stated, an algebraic surface will be understood to mean
a projective algebraic surface over an algebraically
closed field k.
a) Resolution of singularities. The first rigorous
proofs of the existence of a non-singular model were
given in the nineteen-thirties by R. Walker and O. Zar-
iski. The previous geometrical proofs given by the
Italian workers had contained lacunae [9]. The proof of
the existence of a non-singular model of an algebraic
surface for the case of a ground field with positive
characteristic was given in 1956 (cf. Resolution of singu-
larities).
Any non-singular algebraic surface can be
birationally projected in CP3 on a surface with ordi-
nary singularities - ordinary double curves with a fin-
ite number of ordinary cusp points and triple points,
which are triplanar singular points of the surface. Such
surfaces represented the principal object of studies dur-
ing the classical stage of development of the theory of
algebraic surfaces. Owing to the theorem on the resolu-
tion of singularities of algebraic surfaces, global
methods of the theory of schemes may be applied to
the local study of singular points of algebraic surfaces.
The fundamental work on this subject was done by D.
Mumford (1961), who deduced important invariants of
two-dimensional singularities and proved that the tOp<r
logical space of a normal algebraic surface is a topolog-
ical manifold if and only if the surface is non-singular.
b) Numerical invariants of algebraic surfaces. The
concept of the geometric genus Pg(V) of an algebraic
surface is a generalization of the concept of the genus
of an algebraic curve in the theory of algebraic sur-
111
ALGEBRAIC SURFACE
faces. Here, Pg( V) is the maximum number of linearly
independent regular two-dimensional differentials on V
(cf. Geometric genus). The virtual arithmetic genus of
the divisors in the canonical system I K y I of an alge-
braic surface V is known as the linear genus and is
denoted by p(l). For a non-singular algebraic surface
the following equation is valid:
p(1) = (~)+ I = deg(cD+ I,
where CI = - Ky is the first Chern class of the tangent
bundle to the surface V. The multiple canonical sys-
tems I Ki I = I iKy I are called pluri-canonical systems
of the algebraic surface V. The number dim I Ki I + 1 is
referred to as the i-th (pluri-) genus and is denoted by
p.. it is identical with the dimension of the space of
" The number i = dim IKy - D I is known as the special-
everywhere-regular i-tuple two-dimensional differential
forms on V.
According to the Cayley- Noether postulation for-
mula, the number N m of linearly independent forms of
a sufficiently high degree m passing through a space
curve of degree d, genus p and with t triple points and
'T double points, is given by the formula
N m = [mj3] -dm+2t+'T+p-1.
The concept of the arithmetic genus Pa(V), defined
by the formula
Pa(V) = [nil] -d(n-I)+2t+'T+p-1
(where d is the degree, 'T is the number of double
points, p is the genus and t is the number of triple
points of the double curve V; for non-singular surfaces
t =r =d=O and p = 1), was introduced in 1875 for alge-
braic surfaces with ordinary singularities. Other
equivalent definitions of pi V) were given subsequently
(cf. Arithmetic genus).
The difference Pg(V)-Pa(V) is always non-negative;
it is called the irregularity of the algebraic surface V and
is denoted by q(V). An algebraic surface is called regu-
lar if q(V)=O and irregular if q(VO. The first exam-
ples of irregular surfaces date back to the end of the
nineteenth century (A. Cayley, Castelnuovo). Irregular-
ity can be described as the maximal insufficiency of the
linear system which is cut out by the system I C + K v I
on the curve C (Enriques, 1896, see [9]). A
cohomological definition of pa and q for non-singular
algebraic surfaces is given by the formulas:
Pa(V) = -dimk H1(V, ~ v)+di..m;.; H2(V, ~ v)
= xCV, ~ v)-l,
q(V) = dimkHI(V, ~v).
The arithmetic genus Pa(V) of a non-singular alge-
braic surface V is expressed in terms of the Chern
classes of the tangent bundle to V by the formula:
112
deg(c1 )+deg(c2)
I +Pa(V) = 12 '
called Noether's formula. The number I=deg(c2)-4 is
known as the Zeuten - Segre invariant.
c) The Riemann - Roch theorem for algebraic surfaces.
The generalization of the Riemann - Roch theorem for
algebraic curves to algebraic surfaces is due to Castel-
nuovo (1897). For a non-singular algebraic surface V
and a complete linear system I D I of divisors of vir-
tual degree n =(D2) and virtual genus
'1T={(D 2 )+(DKy )} /2+1, the Riemann-Roch
theorem asserts the inequality
dim I D I ;;;;. n-?T+Pa(V)+I-i.
ity index of the divisor D. The cohomological proof of
this theorem (cf. Riemann- Roch theorem) permits the
stipulation of a necessary and sufficient condition such
that this inequality becomes an equality. This condition
is the vanishing of the cohomology group
HI(V, (!) y(D. It has been proved that, for any divisor
D and all sufficiently large numbers n, one has
Hi(V, (!) y(D +nH=O, where i>O and where H is a
hyperplane section of the surface V (J.P. Serre, 1955).
If the characteristic of the ground field is zero, then the
inequality becomes an equality for any ample divisor
D. This result is a generalization of the classical
Picard - Severi theorem concerning the regularity of the
adjoint system [9].
d) Systems of curves on an algebraic surface. The
principal method used by Italian geometers was their
theory of linear and algebraic systems of curves on an
algebraic surface. The concepts of algebraic and linear
equivalence of divisors are closely connected with this
theory. These concepts only coincide for regular alge-
braic surfaces (Castelnuovo, 1896). Any sufficiently
general curve on an algebraic surface V with irregular-
ity q is contained in a maximal algebraic family which
fibrates into linear systems of equal dimension, the
base of the fibration being a variety of dimension q
(Enriques, 1904). This variety is an Abelian variety, and
is known as the Picard variety of the algebraic surface
V (cf. Picard scheme). The proof of Enriques' result
(known as the 'fundamental theorem in the theory of
irregular surfaces', or the 'theorem of completeness of
characteristic series'), as well as subsequent prcofs by
Severi, contain lacunae. The first rigorous transcenden-
= tal proof of this theorem was given by Poincare in
1910, and the algebraic proof came 50 years later (A.
Grothendieck, see [5]). If the characteristic of the
ground field is positive, this theorem is not valid in the
general case (J. Igusa, 1955). It has been shown that, in
the general case, q is not less than the dimension of the
Picard variety.

The non-identity of algebraic and linear equivalence
on an algebraic surface implies that everywhere-regular
one-dimensional differentials exist on it. The dimension
of the space of such differentials in the case where
k=C coincides with the irregularity (Castelnuovo,
Severi, 1905). This fact is a special case of the equality
dimHP(V, gq)=dimHq(V, (l'), which is valid for an
arbitrary compact Kibler manifold V. If char k >0, this
result is invalid.
The group of divisor classes with respect to algebraic
equivalence is finitely generated (see Neron-Severi
group). Its rank is denoted by p and is called the Picard
number of the algebraic surface V. The order of the
torsion subgroup of this group is denoted by (1 and is
known as the Severi divisor of V [9].
Severi also laid the foundations of the theory of
zero-dimensional cycles on an algebraic surface; this
theory contains many interesting problems that have
not yet been solved.
e) Topological properties of an algebraic surface. A
non-singular projective algebraic surface over the field
of complex numbers is a compact four-dimensional
oriented real manifold; in particular, the integral
homology groups Hi(V, Z) and the integral cohomology
groups Hi(V, Z), satisfying Poincare duality, have been
defined for such an algebraic surface. Each divisor on
an algebraic surface V specifies a two-dimensional
cycle, the divisors which are algebraically equivalent to
zero correspond to cycles homologous to zero. The
intersection index of divisors is identical with the topo-
logical intersection index of cycles. The following
equalities are valid: (1 is equal to the order of the tor-
sion group of the group H 2(V, Z) and b 1(V)=2q(V),
where b 1(V) is the first Betti number of the surface V
[4], [9].
A method, based on the study of a linear pencil of
curves on a surface and on the study of the topology of
the curves which vary in this pencil, was developed for
the investigation of the topology of algebraic surfaces.
It was proved in this way that non-singular hyperplanes
in CP3 are simply connected [6]. This method was sub-
sequently applied to studies on the topology of mani-
folds of higher dimension.
The study of etale cohomology of schemes makes it
possible to formulate and prove algebraic analogues of
many classical statements in the topology of algebraic
surfaces for surfaces defined over a field of arbitrary
characteristic. In particular, the l-adic cohomology
groups Hj( V), the fundamental group and the homotopy
type of an algebraic surface have all been defined [9].
E. Picard [6] initiated the theory of integrals on an
algebraic surface. His results were subsequently general-
ized to manifolds of higher dimension (W.V.D. Hodge,
ca. 1940). These results yielded, in particular, the ine-
ALGEBRAIC SURF ACE
quality b 2 (V);;;'p+2pg(V). It was proved, moreover,
that the number of positive squares in the quadratic
form which is defined on H 2(V, Z) by the intersection
index equals 2pg(V) + 1. The analogues of these results
for fields of finite characteristic are not true. The
development of the transcendental theory of algebraic
surfaces is due to the work of P.A. Griffiths.
f) Projective immersions of an algebraiC surface. In
the context of the abstract concept of an algebraic
variety there arose the problem of the existence of
(abstract) algebraic surfaces that are not immersible in
a projective space. It was proved that a non-singular
complete algebraic surface is always projective [8].
There exist complete singular non-projective algebraic
surfaces. There are various numerical criteria for the
projectivity of an algebraic surface (cf. Ample sheaf).
Any non-singular two-dimensional algebraic space is an
algebraic surface.
Minimal models of algebraic surfaces. A non-singular
algebraic surface V is called a minimal model if each
birational morphism V ~ V' onto a non-singular alge-
braic surface V' is an isomorphism. A minimal model
exists in each class of birational equivalence. Except
for classes of ruled surfaces, this minimal model is
unique up to an isomorphism [1], [2], [8]. For an alge-
braic surface to constitute a minimal model it is neces-
sary and sufficient for it not to contain exceptional
curves of the first kind, i.e. irreducible curves which do
contract to a non-singular point for some birational
morphism. Such curves were first studied by M.
Noether in 1895. They are characterized by the follow-
ing requirements: C must be a rational non-singular
curve and (C 2 )= -1 ([1], [2]). Minimal models of ruled
surfaces have been completely classified.
Classification of algebraic surfaces. According to the
classification of Enriques [2], each non-singular alge-
braic surface over a field of characteristic zero belongs
to one of the following types, up to a birational
equivalence: a) ruled surfaces (cf. Ruled surface); b)
two-dimensional Abelian varieties (cf.Abelian variety);
c) K3-surfaces (cf. K3-surface); d) elliptic surfaces (cf.
Elliptic surface); or e) general-type algebraic surfaces
(cf. General-type algebraic surface). This classification
is highly analogous to the classification of algebraic
curves. Moreover, rational curves correspond to ruled
surface, surfaces of the types b) - d) correspond to
elliptic curves, and surfaces of general type correspond
to curves of genus g> 1. The type of the minimal
model of an algebraic surface is determined by the
values of its numerical invariants. Ruled surfaces are
characterized by the condition p 12 = 0; type-b) algebraic
surfaces are characterized by the conditions P12 =1,
Pg = 1 and pa = - 1; type-c). algebraic surfaces are
113
ALGEBRAIC SURFACE
characterized by the conditions Kv=O and q=O; type-
d) algebraic surfaces are characterized by the condi-
tions (Kf.,)=O andpl2>1 orp12=1 andPg=O. Finally,
algebraic surfaces of general type are characterized by
the conditions (Kf.,0 and p 12 > 1.
The class of ruled surfaces comprises rational sur-
faces (cf. Rational surface), which are characterized by
the conditions P2 =Pa =0 (this makes it possible to give
a definite answer to the Liiroth problem in the two-
dimensional case).
Many different results in the theory of the classifica-
tion of algebraic surfaces concern the problem of the
construction of an algebraic surface with given numeri-
cal invariants. Algebraic surfaces are most frequently
constructed by representing them as a double covering
of the projective plane with a specially selected branch
curve (cf. Double plane). The values that can be
assumed by the invariants p(l) and Pa of a general-type
surface are not known (1986). Some of the results
obtained in connection with the classification have
recently been extended to include fields of arbitrary
characteristic. K. Kodaira generalized the classification
results obtained by Enriques to include complex-
analytic surfaces (cf. Analytic surface (in algebraic
geometry.
The moduli problem for algebraic surfaces. This is the
problem of the classification of algebraic surfaces up to
an isomorphism. The first formula for the number of
parameters (moduli) which determine an algebraic sur-
face with given invariants Pa' Pg and p(l) for one class
of algebraic surfaces was given by Noether in 1888.
Non-ruled surfaces with given invariants Pa' Pg and p(l)
depend on M moduli, where
M = lOpa -Pg-2p(l) + 12+8, (*)
where (J is a birational invariant of the algebraic surface
under consideration [2], [9].
The modem theory of deformations gives the follow-
ing interpretation of this classical result. The number M
of moduli is identical with the dimension of the Zariski
space tangent to the local (or, in the algebraic case, the
formal) moduli scheme S a for the polarized surface
(V, a) of this class of algebraic surfaces. Moreover, if
is the canonical extension of the tangent sheaf Tv to
the surface V by the structure sheaf (!) v, defined by the
fundamental class of the polarization a, then
M = dim Im(Hl(V, Ea)~Hl(V, Tv.
Formula (*) is a consequence of this exact sequence
and the Riemann - Roch theorem. The number (J turns
out to be the sum
dimk H2(V, Ea)+dimk H2(V, Tv).
114
Also, the inequality
~H2(V,Ea);;;' 2pg-Pa-1
is valid [9]. The local moduli scheme S a can be singu-
lar, even if char(k)=O. This shows that the 'actual'
number of moduli, i.e. M' = dim Sa, can be less than
M. The difference w = M - M' is called the number of
obstructions to the deformation; the estimate
wE;;;dimH 2(V, Ea) is known. The existence of a global
moduli variety of an algebraic surface has been proved
only for certain cases. The moduli variety for general-
type surfaces and K 3-surfaces exists as an analytic
space or as an algebraic space.
Automorphisms of algebraic surfaces. The automor-
phism group Aut(V) of a complete algebraic surface V
is the group of k-points of a certain group scheme, the
connected component Auto(V) of which is an algebraic
group. If V is not a ruled surface and if
dim Auto(VO, then Pa = -1. If Pg=l=I or if Pg= I
and P2 > I, then V is an elliptic surface and Auto(V) is
a one-dimensional Abelian variety. In the other cases V
and Auto(V) are Abelian surfaces [2]. For general-type
surfaces Aut(V) is a finite subgroup of a projective
group. The group Aut(V) for K3-surfaces over the field
of complex numbers and for ruled surfaces has been
thoroughly studied. If the surface V is not ruled, then
Aut(V) coincides with the group of birational
transformations of V. This group has no algebraic
structure for ruled surfaces and has not been
thoroughly studied (cf. Cremona group). The automor-
phism group of affine algebraic surfaces (cf. Algebraic
variety, automorphism of an) is now under intensive
study.
Algebraic surfaces over algebraically non-closed fields.
Number-theoretic problems in the theory of algebraic
surfaces involve Diophantine problems (cf. Diophantine
geometry). The classification of rational surfaces over
non-closed fields, begun by Enriques, Comessatti and
Segre, has now been completed. The group of birational
automorphisms has been studied for a certain class of
such stafaces.
The results of the theory of algebraic surfaces are
used in the study of algebraic curves over function
fields (d. MordeU conjecture). The generalization of
certain results of the theory of algebraic surfaces
(minimal models, intersection theory) to the wider class
of regular two-dimensional schemes [7] permits the use
of the geometrical language in the study of algebraic
curves over number fields.
References
[I] 'Algebraic surfaces', Trudy Mat. Inst. Steklov. 75 (1965) (in
Russian).
[2] ENRIQUES, F.: 1 superjicie algebraiche, Bologna, 1949.
[3] lUNG, H.W.E.: Aigebraische Fiiichen, Hannover, 1925.

[4] LEFSCHETZ, S.: L'anolysis situs et la geometrie a1gebrique, Paris,
1924.
[5] MUMFORD, D.: Lectures on curves on an algebraic surface,
Princeton Univ. Press, 1966.
[6) PICARD, E. and SIMART, G.: Theorie desfonctions algebriques
de deux variables indepentiantes, 1-2, Chelsea, reprint, 1971.
[7] SHAFAREVICH, I.R.: Lectures on minimal models and birational
transformtJtions of two-dimensional schemes, Tata Inst. Fundam.
Res., 1966.
[8) ZAIusIa, 0.: Introduction to the problem of minimal models in
the theory of algebraic surfaces, Math. Soc. Japan, 1958.
[9) ZAIuSKI, 0.: Algebraic surfaces, Springer, 1971.
[10] ZAIuSKI, 0.: An introduction to the theory of algebraiC surfaces,
Springer, 1969.
1. V. Dolgachev
Editorial comments. One of the more important achieve-
ments is the proof of the Miyaoka- Yao- Bogomolov ine-
quality ~ o;;;;3C2, where Cj denotes the i-th Chern class of an
algebraic surface [A1]. A recent reference is [A2].
References
[A1] MIYAOKA, Y.: 'On the Chern numbers of surfaces of general
type', Invent. Math. 42 (1977), 225-237.
[A2] BAllTH, W., PETERs, C. and YEN, A. VAN DE: Compact com-
plex surfaces, Springer, 1984.
[A3] BBAUVll.LB, A.: 'Surfaces algebrique complexes', Asterique
54 (1958). Translated into English as a London Math. Soc.
Lecture Notes.
[A4] GllIFFITIIS, P.A. and HAIuus, J.E.: Principles of algebraic
geometry, 1-2, Wiley, 1978.
AMS 1980 Subject Classification: 14JXX
ALGEBRAIC SYSTEM - A set with operations and
relations defined on it. An algebraic system is one of
the basic mathematical concepts and its general theory
has been developed in depth. This was done in the
nineteen-fifties, and the work took place on the inter-
face between algebra and mathematical logic.
Fundamental concepts. An algebraic system is an object
A= <A, 0, R> consisting of a non-empty set A, a
family 0 of algebraic operations (cf. Algebraic opera-
tion) oj:An'..",A (iEI) and a family R of relations (cf.
Relation) rj ~A (j EJ) defined on A. The indices
mj
nj, mj of the Cartesian powers of A under consideration
are assumed to be non-negative integers and are said to
be the arities of the respective operations and relations.
The set A is called the carrier or the underlying set of
the algebraic system A, while its elements are called the
elements of this system. The cardinality I A I of A is
said to be the cardinality or order of the algebraic sys-
tem A. The image oj(al,'" ,an) of the element
(aJ, ... ,an,)EA n, under the mapping oj:An'..",A is
called the value of the operation OJ at the point
(aJ, ... ,an,). Similarly, if (aI, ... ,am)Erj' then one
says that the elements a I, . . . ,a of A are in relation
mj
rj' and one writes ria J, ,am). The operations OJ
(i EI) and the relations rj (j EJ), unlike other opera-
tions and relations that may be defined on A, are called
basic or primitive.
ALGEBRAIC SYSTEM
The pair of families <{nj: iE/}; {m/ jEJ}> is
called the type of the algebraiC system A. Two algebraic
systems A, A' have the same type if 1=]', J =f and
nj =n; , mj =mj for all i EI, j EJ. Basic operations OJ,
0; and basic relations rj' rj of two algebraic systems A,
A' of the same type, with identical indices in I, J
respectively, are called similar.
An algebraic system A is called finite if the set A is
finite and it is called of finite type if the set I U J is
finite. An algebraic system A of finite type is written as
A = <A; OJ, ,os, rJ, ,rt >.
An algebraic system A= <A, 0, R> is called a
universal algebra or an algebra if the set R of its basic
relations is empty, and it is called a model or a rela-
tional system if the set 0 of basic operations is empty.
Classical algebraic systems are groups, rings, linear
spaces, linear algebras, totally ordered sets, totally
ordered groups, lattices, etc.
A non-empty subset B of the underlying set A of an
algebraic system A = <A, 0, R> is called closed if for
any elements b J, . ,bn, of B the value oj(b J, . ,bn,}
of each basic operation OJ EO also belongs to B. By
considering the operations of 0 and the relations of R
on a closed subset B, one obtains an algebraic system
B= <B, 0', R' > which is of the same type as the
given algebraic system and is called a subsystem of the
algebraic system A. Subsystems of algebras are called
subalgebras, while subsystems of models are called sub-
models. The concept of a subalgebra strongly depends
on the set of basic operations of the algebra under con-
sideration. Thus, a groupoid G is an algebra of type
<2>, i.e. an algebra with one basic operation
G X G..",G. A groupoid H with a distinguished unit e is
an algebra of type <2,0>, the distinguished element
of which has the property o(e, h)=o(h, e)=h for all
h EH with respect to the basic operation 0: H X H..",H.
For this reason any subgroupoid of a groupoid H with
a distinguished unit e contains e, while a subgroupoid
of a groupoid <H, 0> does not necessarily contain e.
Unlike algebras, any non-empty subset of a model may
be considered as a submodel.
An algebraic system A is isomorphic to an algebraic
system A' of the same type if there exists a one-ta-one
mapping cp of the set A into the set A' such that
q,(o;(at. ... ,an,) = O;(q,(al),' .. ,q,(an,), (1)
riat. ... ,am) <=> rj(q,(ad, . .. ,q,(am)), (2)
holds for all a I, a2, ... , of A and for all i EI, j EJ. A
mapping cp possessing these properties is called an iso-
morphism.
In the following, a class of algebraic systems will be
understood to mean only an abstract class, i.e. a class
of algebraic systems of the same type which, whenever
115
ALGEBRAIC SYSTEM
it contains a system A, also contains all systems iso-
morphic to A. In considering a class Sf of algebraic sys-
tems, all systems in this class are usually written in a
definite signature, as follows. Let the type of the class
Sf be <{nj: iEI}; {m/ JEJ}>. To each iEI is
assigned some symbol Fj, called a functional, while to
each j EJ is assigned a symbol Pj' called a predicate. If
an algebraic system A belongs to the class st and if
OJ: A II, ~A is a basic operation in it, then the element
oj(a], ... ,all) of A is written as F;(a], ... ,all). Simi-
larly, if rj CA mj is a basic relation in A and the element
(a], ... ,am,)Erj' then one writes P/a], ... ,am)=T
(true) or simply P/a], ... ,am). If, on the other hand,
(a], ... ,am)flrj' one writes P/a], ... ,am)=F(false)
or -,P/a] , ... ,am). Let Qj = {Fj: i EI},
Qp = {Pj : j EJ} and let v be the mapping of the union
Qj U Qp into the set of natural numbers {a, 1,2, ... },
defined by the formulas v(Fi)=ni (iEI), v(Pj)=mj
(j EJ). The object Q= <Qj , Qp' v> is the signature of
the class st. A finite signature is written as a string
< 1'1
---'II,) ---'II,). p(m,)
, ... ,1'~ , ] , ... , p(m,)
lo > r , more bne
<F], ... ,Fs; PI, ... ,PI>. An algebraic system A,
written in the signature Q, is known as an Q-system and
is denoted by A= <A, Q>.
The conditions (I) and (2) for an isomorphism of the
systems A and A' become simplified if the systems are
considered in one signature Q. Thus, if Fi (i EJ) and Pj
(j EJ) are the signature symbols, equations (1) and (2)
assume the form
cp(Fj(a j, ... ,an)) = Fj(cp(a j), ... ,CP(an)),
Pia], ... ,am) ~ Picp(aj), ... ,CP(am)).
A homomorphism of an Q-system A into an Q-system
A' is any mapping 4>: A ~A' which fulfils both condi-
tion (3) and the condition
Pia], ... ,am) ~ Picp(aj), ... ,cp(amj ))
for all Fi EQj, Pj EQp and for all a I, a2, ... , of A. A
homomorphism 4>: A~A' is called strong if, for any ele-
ments b I, . . . ,bm, of A' and for any predicate symbol
Pj of Qp, the relation P/ b J, . . . ,bm) = T implies the
existence in A of pre-images a], ... , am, of the ele-
ments b l , . . . ,bm, such that Pj(a], . .. ,am)=T.
, The
concepts of a homomorphism and of a strong
homomorphism of algebras are identical. Models have
homomorphisms that are not strong, and have one-to-
one homomorphisms that are not isomorphisms. Under
a homomorphism 4>: A~A', images of subsystems of A
are subsystems in A'; and non-empty complete pre-
images of subsystems of A' are subsystems in A.
An equivalence relation () CA X A is called a
congruence of the Q-system A if
116
8(a], bd, . ; . ,8(an" bn) ~
~ 8(Fj(a), ... ,an,),Fj(b), ... ,bn,))
for all a J, b], ... ,all" bll, of A and all F j EQj . For
each homomorphism 4> of an algebraic system A the
binary relation fJ(a, b) which holds if and only if
4>(a)=4>(b) is a congruence in A, called its kernel
congruence. For any congruence 0 of an Q-system A
and for each element a EA the set
a / fJ = {x EA: fJ(x, a)} is called a coset of the algebraic
system A by the congruence (). On the assumption that,
for each F; EQj' Pj EQp'
Fj(aj /fJ, ... ,an, /fJ) = Fj(a], ... ,an)/fJ
and
P/b j /fJ, ... ,bmj /fJ) = T
if and only if there exist elements c], ... ,cmj in A such
that fJ(b] , cd, ... ,fJ(bm , cm ) and Pj'(CI, ... ,cm ), one
j j ,
obtains the algebraic system A / fJ, which is of the same
type as the given system; it is called the quotient system
of the algebraic system A by the congruence fJ. For
. fly,
each congruence () of an algebraic system A the canoni-
cal mapping 4>( a) = a / () (a EA) is a homomorphism of
A onto the quotient system A / fJ for which the given
congruence () is the kernel congruence. If 4> is a
homomorphism of an algebraic system A into an alge-
braic system A' and () is the kernel congruence for 4>,
then the mapping 1/;{a / fJ)=4>(a) is a homomorphism of
the quotient system A / fJ into the algebraic system A'.
If the homomorphism 4> is strong, 1/; is an isomorphism.
(3) The Cartesian product of Q-systems Aa = <Aa, Q>,
aEA=I= 0, is the Q-system D= <D, Q> in which D is
(4)
the Cartesian product of the underlying sets A a (a E A)
and the basic operations and the basic relations on D
are given by the conditions: F;(dJ, ... ,dn)
(d l , . . . ,dn, ED, FiEAf) is the element dED with coor-
(5) dinates d(a)=Fi(dl(a), ... ,dn,(a)) (aEA),
P/dJ, ... ,dm)=T (PjEQp) if and only if
P/dl(a), ... , dm,(a)) = Tfor all aEA.
A language of the first order. The basic formal
language of the theory of algebraic systems is the first-
order language L which is constructed as follows. The
alphabet of the language L in the given signature
Q=<g,j,Qp'v>, Qj={Fi: iE/}, g,p={Pj:jEJ}, con-
sists of the object variables x I, X2, . . . , the function
symbols Fi (i EI), the predicate symbols Pj (j EJ), the
logical connectives
the quantifiers:
\;JXk - 'for each element Xk';
3Xk - 'there exists an element Xk';
and the auxiliary symbols: brackets and commas.

The (first-order) properties of n-systems are expressed
by finite sequences of alphabet symbols, or words, con-
structed in accordance with specific rules and known as
terms and formulas. It is inductively assumed that each
word of the form Xk or F; is a term if p(Fj)=O; if
fl> ... ,In are terms and if n =p(Fj), then
Fj(fl> ... ,In) is also a term.
If A is an n-system and if f(x I, . . . ,xn ) is a term of
the signature n containing object variables x I, . . . ,Xk,
then if one replaces x I> ,Xk by some elements
a I, . . . ,ak in A, and executes the operations in A on
the latter elements corresponding to the symbols of f n are true.
forming part of this term, then one obtains an element
f(al> ... ,ak) from A, which is called the value of the
term f(xJ, . .. ,Xk) for XI =aJ, ... ,xk=ak. If <I> is a
homomorphism of an n-system A into an n-system A',
then
qij(aJ, ... ,ad) = j(cp(ad, ... ,cp(ad)
The concept of a formula of the signature n, and of
free and bound object variables in it, is also defined
inductively:
I) If P is any predicate symbol for n or the equality
sign =, m = p(P) or 2 respectively, and if fl' ... ,fm
are arbitrary terms of the signature n, then the word
P(fJ, ... ,fm) is a formula in which all object variables
are free.
2) If ~ is a formula, so is ...,~. Free (bound) object
variables in the formula ...,~ are those and only those
variables which are free (bound) in ~.
3) If ~I' fu are formulas and if the object variables
which form part of both formulas are free in both for-
mulas, then the words
(6)
are also formulas.
Object variables which are free (bound) in at least
one of the formulas ~I' ~2 are called free (bound) in
the formulas (6) as well.
4) If an object variable Xk occurs free in a formula ~,
then the words 0txkm, (3Xkm are again formulas in
which the variable Xk is bound, while all the other
object variables which occur free or bound in the for-
mula ~ remain free or bound in the formulas 0txkm,
m
(3Xk as well.
Let an n-system A and a formula ~ of the signature
n be given. If one now assigns to all free object vari-
ables x I, . . . ,Xk from ~ some values a I, . . . ,ak from
A and if one interprets the function and the predicate
symbols forming part of ~ as the respective basic
operations and basic relations in A, then one obtains a
definite statement which may be true or false. One
accordingly assigns to the formula ~ the value T or F,
for XI =aJ, ... ,xk=ak denoted by ffial, . . . ,ak). If
ALGEBRAIC SYSTEM
is an isomorphic mapping of an n-system A into an
<I>
n-system A', then
ffia J, ... ,ak) = fficp(a 1), ... ,CP(ak)),
for all a J, ,ak from A.
A formula ~ is called closed if it contains no free
object variables. For any closed formula ~ of the signa-
ture n and any n-system A one may speak of ~ being
true or false in A. A set S of closed formulas of a
given signature n is said to be realizable or consistent if
there exists an n-system in which all formulas from S
The compactness theorem or the Godel- Mal'tsev
theorem states: If each finite subset of an infinite set S
of closed formulas of a given signature 0 is realizable,
then the entire set S is realizable as well.
Axiomatizable classes. Let S be some set of closed for-
mulas of a signature O. The class of all O-systems in
which all formulas of S are true will be denoted by KS.
The set Th ~ of all closed formulas of a signature 0
which are true in all O-systems of the given class ~, is
called the elementary theory of ~. In particular, if (A)
is the class of n-systems isomorphic to the given 0-
system A, then Th(A) is said to be the elementary
theory of the n-system A and is simply denoted by
Th A. A class ~ of O-systems is said to be axiomatizable
if ~ = K Th~. A class ~ of O-systems is axiomatizable
if and only if there exists a set S of closed formulas of
the signature 0 such that ~=KS.
Along with the general concept ofaxiomatizability,
axiomatizability using first-order formulas of a special
kind may also be considered. The special formulas of a
given signature which are the most important in alge-
bra are:
Identities - the formulas
(Vxd .. (Vxs)P(jJ, ... ,fm),
where P is some predicate symbol from 0 or the equal-
ity sign =, and fl' ... ,fm are terms of the signature 0
in XI, . . . ,XS.
Quasi-identities - the formulas
(VXl) (Vxs)[P(jJ, ... ,jk)& ... &Q(gJ, ... ,gm)~
~R(hJ, ... ,hn )],
where P, ... , Q, R are certain predicate symbols from
Op or equality signs, while
fI. ,/k,gl> ... ,gm, hl> ... ,hn are terms of the sig-
nature 0 in x I> ,XS
Universal formulas - the formulas 0tx I) ... 0txsm,
where ~ is a formula of the signature 0 not containing
quantifiers.
If a set S of identities (quasi-identities or universal
formulas) of a signature n is given, then the class KS is
117
ALGEBRAIC SYSTEM
known as a variety (quasi-variety or universal class) of
Q-systems.
Birkhoffs theorem: A non-empty class ~ of Q-
systems is a variety if and only if it is closed with
respect to subsystems, Cartesian products and
homomorphic images.
If A= <A, Q> is some n-system, then, by exchang-
ing each function symbol F; from nf for a predicate
symbol F; of arity n; + I (higher by one) and putting,
for elements a J, . . ,ani' a from A
F;(a], ... ,ani' a) ~ Fi(a\> ... ,an) = a,
one obtains a model A = <A, n > for which
n;=np U {F;: F;Enf } Submodels of the model A
are called submodels of the n-system A. For any non-
empty finite subsets Aa CA, n,8 en, the model
Aap = <A a , n,8 > is called a finite depletion of the finite
submodel Aa = <A a, n > of the n-system A. An n-
system A is said to be locally imbeddable in a class ~ of
n-systems if for each finite depletion Aap of any finite
submodel Aa of the n-system A there exists in the class
~ an n-system B (depending on the chosen finite deple-
tion Aap) such that the model Aap = <A a, n,8 > is iso-
morphic to the model BaP=<Ba,n,8> for a suitable
subset Ba CB.
A subclass 53 of a class ~ of n-systems is called
universal (or universally axiomatizable) in ~ if there
exists a set S of universal formulas of the signature n
such that 53 = KS n
Sl'.
The Tarski - Los theorem: A subclass 53 of a class ~
of n-systems is universal in ~ if and only if 53 contains
all systems from ~ which are locally imbeddable in 53.
Filtered products. Let D = ITAa be the Cartesian pro-
duct of the n-systems Aa (a E A, A=r'= 0) and let cI> be
some filter over A. The relation
d=lph ~ {a:EA: d(a:) =h(a:)} Eet> (d, h ED)
is an equivalence relation on the underlying set D of
the n-system D. For each element d ED, let d / cI> be
the coset by this equivalence and let
D / cI> = {d / cI>: d ED}. Putting
F,(d j Iet>, .. ,dn, let = det> ~
{a: F;(dj(a), . .. ,dn,(a=d(a)} E et> (FiEfJj ),
$=>
P/d j Iet>, ,dm} let ~ of; Algebraic systems, variety of.
~ {a: P/dj(a), ... ,dm/a} E et> (P;EfJp ),
one obtains the n-system D/cI>=<D /cI>,n>, which
is called the product of the n-systems Au (aEA)filtered
through the filter cI>. The n-systems Au (aEA) are said
to be the factors of this product. If cI> is. an ultrafilter
over A, the filtered product D / cI> is known as the
ultraproduct of the Q-systems Au (aEA).
The ultraproduct theorem: If D / cI> is the ultraproduct
118
of the n-systems Aa (aEA) and if ffixJ, ... ,xd is an
arbitrary formula of the signature n in which
x I, . . . ,Xk are free object variables, then for any ele-
ments d J, . . . ,dk ED:
'is{d j let>, ... ,dk let = T ~
~ {a: 'is{dj(a), ... ,dk(a=T} E et>.
In particular, a closed formula iY of the signature n
is true in the ultraproduct D / cI> of the n-systems Au
(aEA) if and only if the set of numbers of factors in
which the formula iY is true belongs to the ultrafilter cI>.
Therefore, any axiomatizable class of n-systems is
closed with respect to ultraproducts.
A class L of n-systems is universally axiomatizable if
and only if it is closed with respect to subsystems and
ultraproducts.
An n-system E= < {e}, Q> is said to be a unit if its
underlying set consists of a single element, e.g. an ele-
ment e, and if Pie, ... ,e)=T for all PjEnp-
Mal'tsev's theorem: A class ~ of n-systems is a
quasi-variety if and only if it contains a unit n-system
and is closed with respect to subsystems and filtered
products (through an arbitrary filter).
Completeness and categoricity. A non-empty class ~ of
n-systems is called categorical if all n-systems from ~
are mutually isomorphic. Any categorical axiomatizable
class of Q-systems consists of one (up to an isomor-
phism) finite n-system.
A class ~ of n-systems is called categorical in cardi-
nality m if it contains an n-system of cardinality m and
if all Q-systems from ~ of cardinality m are mutually
isomorphic. For example, the class of algebraically
closed fields of fixed characteristic is categorical in any
uncountable infinite cardinality. A non-empty class ~
of n-systems is called complete if for any n-systems
A, B from ~ the equality ThA=ThB is valid.
Vaught's theorem. If an axiomatizable class Sl' of n-
systems is categorical in some cardinality
m;;;;' I Qf U np I and if all n-systems in Sl' are infinite,
then ~ is a complete class.
In particular, the class of all algebraically closed
fields of fixed characteristic is complete.
See also Algebraic system, automorphism of an; Alge-
braic systems, quasi-variety of; Algebraic systems, class
References
[I] MAL'TSEV, A.I.: Algebraic systems, Springer, 1973 (translated
from the Russian).
[2] COHN, P.M.: Universal algebra, Reidel, 1981.
[3] GRATZER, G.: Universal algebra, Springer, 1979.
[4] BELL, J.L. and SLOMSON, A.B.: Models and ultraproducts: an
introduction, N orth- Holland, 1971.
[5] CHANG, c.c. and KEISLER, H.J.: Model theory, North-Holland,
1973.
D.M. Smirnov

Editorial comments. The terminology in use for this sub-
ject in English-speaking countries differs in a number of
respects from the Russian terminology. What is here called
an algebraic system is usually called a (first-order) structure
(a structure of type 0, if a given signature or similarity type
o has been specified); the term algebra is often used
instead of 'structure' in the case when 0 has no primitive
relations. The term model is reserved for a structure which
satisfies a given set of first-order sentences (closed formu-
las); it is not normally used in the sense defined above. The
underlying set of a structure is sometimes called its
universe.
Although the text is correct in saying that the subject was
largely developed in the nineteen-fifties (principally through
work of LA Henkin [L.A Khenkin], AI. Mal'tsev, A Robin-
son and A. Tarski), two important precursors dating from
1935 should be mentioned. They are Birkhoff's theorem
characterizing varieties of algebras [A1] and Tarski's defini-
tion of validity of formulas in a first-order structure [A2].
References
[A 1] BIRKHOFF, G.: 'On the structure of abstract algebras', Proc.
Cambridge Phi/os. Soc. 31 (1935), 433-454.
[A2] TARSKI, A.: 'Der Wahrheitsbegriff in den formalisierten Spra-
chen', Studia Phi/os. 1 (1935),261-405.
AMS 1980 Subject Classification: 08AXX, 03CXX
ALGEBIlAIC SYSTEM, AUTOMORPIDSM OF AN -
An isomorphic mapping of an algebraic system onto
itself. An automorphism of an O-system A= <A, 0> is
a one-to-one mapping cp of the set A onto itself having
the following properties:
CP(F(XIo . .. ,xn = F(CP(XI), ... ,cp(xn , (1)
P(XIo . .. ,xm ) ~ P(CP(Xl), ... ,q,(xm )),
for all x], x2, ... , from A and for all F, P from O. In
other words, an automorphism of an O-system A is an
isomorphic mapping of the system A onto itself. Let G
be the set of all automorphisms of the system A. If
cpEG, the inverse mapping cp-I also has the properties
(1) and (2), and for this reason cp-I EG. The product
a=# of two automorphisms cp, t/I of the system A,
defined by the formula a(x)=t/I(cp(x, xEA, is again an
automorphism of the system A. Since multiplication of
mappings is associative, < G, " -I > is a group, known
as the group of all automorphisms of the system A; it is
denoted by Aut(A). The subgroups of the group Aut(A)
are simply called automorphism groups of the system A.
Let cp be an automorphism of the system A and let 8
be a congruence of this system. Putting
(x,Y)E8.p ~ (cp-l(X), cp-I(y))E8, x,YEA,
one again obtains a congruence 8.p of the system A. The
automorphism cp is known as an Ie-automorphism if
8.p = () for any congruence () of the system A. The set
IC(A) of all IC-automorphisms of the system A is a
normal subgroup of the group Aut(A), and the quotient
ALGEBRAIC SYSTEM, AUTOMORPHISM OF AN
group Aut(A) / IC(A) is isomorphic to an automor-
phism group of the lattice of all congruences of the sys-
tem A [1]. In particular, any inner automorphism
x ~a -I xa of a group defined by a fixed element a of
this group is an IC-automorphism. However, the exam-
ple of a cyclic group of prime order shows that not all
IC-automorphisms of a group are inner.
Let ~ be a non-trivial" variety of O-systems or any
other class of O-systems comprising free systems of any
(non-zero) rank. An automorphism cp of a system A of
the class ~ is called an I -automorphism if there exists a
term f.p(xlo'" ,xn ) of the signature 0, in the
unknowns XI, . ,xn , for which: 1) in the system A
there exist elements a 2, . . . ,an such that for each ele-
ment x EA the equality
q,(x) = !.p(x,a2,'" ,an)
is valid; and 2) for any system B of the class ~ the
mapping
is an automorphism of this system for any arbitrary
selection of elements X2, . . . ,Xn in the system B. The
set I(A) of all I-automorphisms for each system A of
the class ~ is a normal subgroup of the group Aut(A).
In the class ~ of all groups the concept of an 1-
automorphism coincides with the concept of an inner
automorphism of the group [2]. For the more general
concept of a formula automorphism of O-systems, see [3].
Let A be an algebraic system. By replacing each
basic operation F in A by the predicate
R(XIo .. . ,xn,y) ~ F(XIo ... ,xn)=y
(2)
(X., ... ,xn,YEA),
one obtains the so-called model A* which represents the
system A. The equality Aut(A*)=Aut(A) is valid. If
the systems A= <A, 0> and A' = <A, 0' > have a
common carrier A, and if 0 CO', then Aut(A) d Aut(A').
If the O-system A with a finite number of generators is
finitely approximable, the group Aut(A) is also finitely
approximable (cf. [1 D. Let ~ be a class of O-systems
and let Aut(~) be the class of all isomorphic copies of
the groups Aut(A), AE~, and let SAut(~) be the class
of subgroups of groups from the class Aut(~). The class
SAut(~) consists of groups which are isomorphically
imbeddable into the groups Aut(A), AE~.
The following two problems arose in the study of
automorphism groups of algebraic systems.
1) Given a class ~ of O-systems, what can one say
about the classes Aut(~) and SAut(~)?
2) Let an (abstract) class K of groups be given. Does
there exist a class ~ of O-systems with a given signature
o such that K=Aut(~) or even K=SAut(~)? It has
been proved that for any axiomatizable class ~ of
models the class of groups SAut(~) is universally
119
ALGEBRAIC SYSTEM, AUTOMORPHISM OF AN
axiomatizable [1]. It has also been proved [1], [4] that if
~ is an axiomatizable class of models comprising infin-
ite models, if <B, OS;;; > is a totally ordered set and if G
is an automorphism group of the model <B, OS;;; >,
then there exists a model AE~ such that A ~B, and for
each element g E G there exists an automorphism cp of
the system A such that g(x)=cp(x) for all XEB. The
group G is called 1) universal if GESAut(~) for any
axiomatizable class ~ of models comprising infinite
models; and 2) a group of ordered automorphisms of an
ordered group H (cf. Totally ordered group) if G is iso-
morphic to some automorphism group of the group H
which preserves the given total order OS;;; of this group
(i.e. a os;;;b=>cp(a)os;;;CP(b) for all a, bEH, CPEG).
Let I be the class of totally ordered sets <M, OS;;; >,
let U be the class of universal groups, let RO be the
class of right-ordered groups and let OA be the class of
ordered automorphism groups of free Abelian groups.
Then [4], [51 [6]:
SAut(l) = U = RO = ~A.
Each group is isomorphic to the group of all automor-
phisms of some ~-algebra. If ~ is the class of all rings,
Aut(~) is the class of all groups [1]. However, if ~ is
the class of all groups, Aut(~)=F~; for example, the
cyclic groups C3 , C5 , C7 of the respective orders 3, 5
and 7 do not belong to the class Aut(~). There is also
no topological group whose group of all topological
automorphisms is isomorphic to Cs [7].
References
[1] PLOTKIN, B.I.: Groups of automorphisms of algebraic systems,
Wolters-Noordhoff, 1972.
[2] CsAKANY, B.: 'Inner automorphisms of universal algebras',
Pub!. Math. Debrecen 12 (1965),331-333.
[3] GRANT, J.: 'Automorphisms definable by formulas', Pacific J.
Math. 44 (1973), 107-115.
[4] RABIN, M.O.: 'Universal groups of automorphisms of models',
in Theory of models, North-Holland, 1965, pp. 274-284.
[5] COHN, P.M.: 'Groups of order automorphisms of ordered sets',
Mathematika 4 (1957), 41-50.
[6] SMlRNOV, D.M.: 'Right-ordered groups', Algebra i Logika 5,
no. 6 (1966), 41-59 (in Russian).
[7] WILLE, R.J.: 'The existence of a topological group with auto-
morphism group C 7 ', Quart. J. Math. Oxford (2) 18 (1967),
53-57. D. MS'
. mlrnov
AMS 1980 Subject Classification: 08A35
ALGEBRAIC SYSTEMS, CLASS OF - A class of
algebraic systems of the same type. All systems of a
given type are assumed to be written in a given signa-
ture ~ and are called ~-systems. A class ~ of 12-systems
is called abstract if, whenever it contains a system A, it
also contains all 12-systems isomorphic to A.
Let ~ be an abstract class of ~-systems. One says
that an 12-system A has a local set of ~-subsystems if
there exists an inclusion-directed set {A,,: ilEA} of
subsystems A" of the system A which cover A (i.e.
120
U a Aa=A) and which belong to~. A class ~ is called
local if each ~-system A with a local set of ~
subsystems belongs to the class ~. Theorems which
establish the local nature of given abstract classes are
called local (cf. MaI'tsev local theorems).
An ~-system A is called ~-approximable (or ~
residual) if, for any predicate P E {~p, =} (i.e. for any
basic predicate as well as for the predicate coinciding
with the equality relation in A) and for any elements
aJ, ... ,an in A for which P(aJ, ... ,an)=F, there
exists a homomorphism cp: A~B of the system A into
some system B of the class ~ for which, again,
P(cp(a)), ... ,CP(an=F. Any subsystem of a ~
approximable system is itself ~-approximable. If ~ is
the class of all finite ~-systems, a ~-approximable sys-
tem 2{ is called finitely apprOXimable (or residually fin-
ite). If an abstract class ~ has a unit system
E = < {e }, ~ >, an ~-system A is ~-approximable if
and only if it is isomorphically imbeddable in a Carte-
sian product of systems from the class ~ [3]. A class ~
is called residual if all ~-approximable systems belong
to the class ~. A class ~ is homomorphically closed if it
contains, for each of its ~-systems A, also all ~-systems
that are homomorphic images of A. All residual
homomorphically-closed classes are local [5].
A class ~ of ~-systems is called (finitely) axiomatiz-
able if there exists a (finite) set S of first-order closed
formulas of the signature 12 such that ~ consists of
exactly those ~-systems in which all formulas of S are
true. Finitely-axiomatizable classes are also referred to
as elementary classes. It has been shown [5] with the aid
of the generalized continuum hypothesis that: 1) a class
~ of algebraic systems is axiomatizable if and only if it
is closed with respect to ultra-products, and its comple-
ment (in the class of all ~-systems) is closed with
respect to ultra-powers; 2) a class ~ of algebraic sys-
tems is elementary if and only if both it and its com-
plement are closed with respect to ultra-products. The
theory ofaxiomatizable classes of algebraic systems
deals with the connection between the structural pro-
perties of these classes and the syntactic features of the
formal language in which these classes may be
specified. Axiomatizable classes which play an espe-
cially important role in algebra include varieties (cf.
Algebraic systems, variety of) and quasi-varieties (cf.
Algebraic systems, quasi-variety of), which are local and
residual.
In addition to axiomatizability by first-order closed
formulas, axiomatizability by special second-order
closed formulas is also considered. To the function and
predicate signature symbols F; (i E/), Pj (j EI) of a
fixed signature 12 are added predicate variables
R 1 , R 2, . . . . Let 0: be a quantifier-free formula of the
first order, consisting of function and predicate signa-

ture symbols, predicate variables R I, . . . , Rs and
object variables x I, . . . , x,. A second-order formula
Q 0:, where Q is some sequence of quantifiers of the
type ('VR i ), (3Ri) or ('Vxd is known as crypto-universal.
The second-order formulas formed from the crypto-
universal formulas without free object variables, with
the aid of the logical connectives &, V,~, -, and with
subsequent quantification by V on all free predicate
variables encountered in the representations of the
crypto-universal formulas, are called Boolean-universal
formulas of the signature n. A class Sl' of n-systems is
called quasi-universal if there exists a set S of Boolean-
universal formulas of the signature n such that Sl' con-
sists of the n-systems in which all formulas of S are
true, and only of such systems. A quasi-universal class
of Q-systems is local (Mal'tsev's theorem). A.I. Mal'tsev
[4] gave a more detailed definition of a quasi-universal
class.
References
[I] MA!.'rsEv, A.I.: 'A general method for obtaining local
theorems in group theory', Uchen. Zap. Ivanovsk. Gas. Ped.
Inst. 1, no. 1 (1941), 3-9 (in Russian).
[2] MAL'rsEv, A.I.: 'Model correspondences', Izv. Akad. Nauk
SSSR Ser. Mat. 23 (1959),313-336 (in Russian).
[3] MA!.'rsEv, A.I.: Algebraic systems, Springer, 1973 (translated
from the Russian).
[4] MAL'rsEv, A.I.: 'Some problems in the theory of classes of
models', in Proc. 4-th AI/-Union Math. Congress 1961, Len-
ingrad, 1963 (in Russian). Trans!. in: Amer. Math. Soc.
Trans!. (2) 83 (1969), 1-48.
[5] COHN, P.M.: Universal algebra, Reidel, 1981.
[6] CLEAVE, J.P.: 'Local properties of systems', J. London Math.
Soc. 44 (1969), 121-130.
D.M. Smirnov
Editorial comments. The characterization (under the
generalized continuum hypothesis) ofaxiomatizable and ele-
mentary classes is due to H.J. Keisler [A1].
Translations of the articles [1], [2] and [4] may (also) be
found in [A2], Chapts. 2, 11 and 26, respectively.
The term inductive class is sometimes used instead of
'local class'.
References
[A1] KEISLER, H.J.: 'Ultraproducts and elementary classes', Indag.
Math. 23 (1961), 477-495.
[A2] MAL'CEV, A.1. [A.1. MAL'rsEv]: The metamathematics of alge-
braic systems. Collected papers: 1936 - 1967, North-Holland,
1971.
AMS 1980 Subject Classification: 08C10, 03C52
ALGEBRAIC SYSTEMS, QUASI-VARIETY OF - A
class of algebraic systems (Q-systems) axiomatized by
special formulas of a first-order logical language, called
quasi-identities or conditional identities, of the form
(\iXl)' .. (VxS>
[P1((fl), ... ,f,,!!)&'" & Pk((fk) , ... ,f'!:') ~
~ Po((f) , ... ,f"?~)],
where Po, ... ,PkEQp U{=}, and frO), ... ,f'!:')
ALGEBRAIC SYSTEMS, QUASI-VARIETY OF
terms of the signature Q in the object variables
x I, . . . , Xs' By virtue of Mal'tsev's theorem [1] a
quasi-variety Sl' of algebraic systems of signature n can
also be defined as an abstract class of Q-systems con-
taining the unit Q-system E, and which is closed with
respect to subsystems and filtered products [1], [2]. An
axiomatizable class of n-systems is a quasi-variety if
and only if it contains the unit Q-system E and is
closed with respect to subsystems and Cartesian pro-
ducts. If Sl' is a quasi-variety of signature n, the sub-
class Sl'l of systems of Sl' that are isomorphically
imbeddable into suitable systems of some quasi-variety
Sl" with signature Q':2 Q, is itself a quasi-variety. Thus,
the class of semi-groups imbeddable into groups is a
quasi-variety; the class of associative rings without zero
divisors imbeddable into associative skew-fields is also
a quasi-variety.
A quasi-variety Sl' of signature n is called finitely
definable (or, is said to have a finite basis of quasi-
identities) if there exists a finite set S of quasi-identities
of Q such that Sl' consists of only those Q-systems in
which all the formulas from the set S are true. For
instance, the quasi-variety of all semi-groups ,vith can-
cellation is defined by the two quasi-identities
zx=zy ~ x=y, xz=yz ~ x=y,
and is therefore finitely definable. On the other hand,
the quasi-variety of semi-groups imbeddable into
groups has no finite basis of quasi-identities [l], [2].
Let Sl' be an arbitrary (not necessarily abstract) class
of Q-systems; the smallest quasi-variety containing Sl' is
said to be the implicative closure of the class Sl'. It con-
sists of subsystems of isomorphic copies of filtered pro-
ducts of Q-systems of the class Sl' U {E}, where E is
the unit Q-system. If Sl' is the implicative closure of a
class ill: of n-systems, ill: is called a generating class of
the quasi-variety Sl'. A quasi-variety Sl' is generated by
one system if and only if for any two systems A, B of Sl'
there exists in the class Sl' a system C containing sub-
systems isomorphic to A and B [l]. Any quasi-variety Sl'
containing systems other than one-element systems has
free systems of any rank, which are at the same time
free systems in the equational closure of the class Sl'.
The quasi-varieties of n-systems contained in some
given quasi-variety Sl' of signature n constitute a com-
plete lattice with respect to set-theoretic inclusion. The
atoms of the lattice of all quasi-varieties of signature Q
are called minimal quasi-varieties of n. A minimal
quasi-variety we is generated by anyone of its non-unit
systems. Every quasi-variety with a non-unit system
contains at least one minimal quasi-variety. If Sl' is a
quasi-variety of Q-systems of finite signature n, all its
sub-quasi-varieties constitute a groupoid with respect to
are the Mal'tsev Sl'-multiplication [3].
121
ALGEBRAIC SYSTEMS, QUASI-VARIETY OF
References
[1] MAL'TSEV,A.I.: Algebraic systems, Springer, 1973 (translated
from the Russian).
[2] COHN, P.M.: Universal algebra, Reidel, 1981.
[3] MAL'TSEV, AI.: 'Multiplication of classes of algebraic systems',
Siberian Math. J. 8, no. 2 (1967), 254-267. (Sibirsk Mat. Zh. 8,
no. 2 (1967), 346-365)
D.M. Smirnov
Editorial comments. In the Western literature, quasi-
identities are commonly called Horn sentences (cf. [A1]).
For a categorical treatment of quasi-varieties, see [A3]; for
their finitary analogue, see [A2]. Mal'tsev's article [A3] may
also be found in [A4] as Chapt. 32.
References
[A1] HORN, A: 'On sentences which are true of direct unions of
algebras', J. Symbolic Logic 16 (1951), 14-21.
[A2] ISBELL, l.R.: 'General functional semantics, I', Amer. J. Math.
94 (1972),535-596.
[A3] KEANE, 0.: 'Abstract Horn theories', in F.w. Lawvere, C.
Maurer and C. Wraith (eds.): Model theory and topoi, Lec!.
Notes in Math., Vol. 445, Springer, 1975, pp. 15-50.
[A4] MAL'CEV, AI. [AI. MAL'TSEV): The metamathematics of alge-
braic systems. Collected papers: 1936 - 1967, North-Holland,
1971.
AMS 1980 Subject Classification: 08C15
ALGEBRAIC SYSTEMS, VARIEI'Y OF - A class of
algebraic systems (cf. Algebraic systems, class of) of a
fixed signature n, axiomatizable by identities, i.e. by
formulas of the type
(Vx I) ... (VXs)P(jh ... ,f,,),
where P is some predicate symbol from n or the equal-
ity sign, while f\, ... ,f,. are terms of the signature n
in the object variables x \, ... ,Xs A variety of alge-
braic systems is also known as an equational class, or a
primitive class. A variety of signature n can also be
defined (Birkhoffs theorem) as a non-empty class of n-
systems closed with respect to subsystems,
homomorphic images and Cartesian products.
The intersection of all varieties of signature n which
contain a given (not necessarily abstract) class ~ of n-
systems is called the equational closure of the class ~
(or the variety generated by the class ~), and is
denoted by var~. In particular, if the class ~ consists
of a single n-system A, its equational closure is denoted
by var A. If the system A is finite, all finitely-generated
systems in the variety var A are also finite [1], [2].
Let !l' be a class of n-systems, let S!l' be the class
of subsystems of systems of !l', let H!l' be the class of
homomorphic images of the systems from !l', and let
II!l' be the class of isomorphic copies of Cartesian pro-
ducts of the systems of !l'. The following relation [1],
[2] is valid for an arbitrary non-empty class ~ of n-
systems:
var st = HS IIR
A variety is said to be trivial if the identity x = y is true
in each one of its systems. Any non-trivial variety we
122
has free systems FmOIR) of any rank m and
[n = var F Ko (IDe) [1], [2].
Let S be a set of identities of the signature n and let
KS be the class of all n-systems in which all the identi-
ties of S are true. If the equality IDe = KS is satisfied for
a variety IDe of signature n, S is known as a basis for
[n. A variety IDe is known as finitely baseable if it has a
finite basis S. For any system A, a basis of the variety
var A is also known as a basis of identities of the system
A. If [n is a finitely-baseable variety of algebras of a
finite signature and if all algebras of IDe have distribu-
tive congruence lattices, then each finite algebra A of
[n has a finite basis of identities [10]. In particular, any
finite lattice <A, V, 1\ > has a finite basis of identi-
ties. Any finite group has a finite basis of identities [3].
On the other hand, there exists a six-element semi-
group [5] and a three-element groupoid [6] without a
finite basis of identities.
The varieties of n-systems contained in some fixed
variety IDe of signature n constitute under inclusion a
complete lattice L(IDe) with a zero and a unit, known as
the lattice of subvarieties of the variety [)C. The zero of
this lattice is the variety with the basis x =y,
P(x\, ... ,xn ) (PEn), while its unit is the variety IDe. If
the variety IDe is non-trivial, the lattice L(IDe) is anti-
isomorphic to the lattice of all fully-characteristic
congruences (cf. Fully-characteristic congruence) of the
system F ~o (IDe) of countable rank which is free in IDe
[1]. The lattice Lo of all varieties of signature n is
infinite, except for the case when the set n is finite and
consists of predicate symbols only. The exact value of
the cardinality of the infinite lattice Lo is known [1].
The lattice of all lattice varieties is distributive and has
the cardinality of the continuum [7], [8]. The lattice of
all group varieties is modular, but it is not distributive
[3], [4]. The lattice of varieties of commutative semi-
groups is not modular [9].
Atoms of the lattice Lo of all varieties of signature n
are known as minimal varieties of signature n. Every
variety with a non-unit system contains at least one
minimal variety. If the n-system A is finite and of finite
type, then the variety var A contains only a finite
number of minimal subvarieties [1].
Let 21:, IS be subvarieties of a fixed variety we of n-
systems. The Mal'tsev product 2!9J(01S denotes the class
of those systems A of we with a congruence () such that
(A/()EIS, and all cosets a/() (aEA), which are sys-
tems in we, belong to 2!. If we is the variety of all
groups and if 21: and IS are subvarieties of it, then the
product 21: IJJ/oSB is identical with the Neumann product
[3]. The product of varieties of semi-groups need not be
a variety. A variety we of n-systems is called polarized
if there exists a term e(x) of the signature n such that
in each system from we the identities e (x) =e (y),

F(e(x), ... ,e(x=e(x) (FED) are true. If ID1 is a.
polarized variety of algebras and the congruences in all
algebras in !In are commutable, then the Mal'tsev pro-
duct 2(ID/0!8 of any subvarieties 2( and !8 in !In is a
variety. One may speak, in particular, of the groupoid
G[(!In) of subvarieties of an arbitrary variety !In of
groups, dngs, etc. If ID1 is the variety of all groups or
all Lie algebras over a fixed field P of characteristic
zero, then G[(!In) is a free semi-group [I].
References
[I] MAL'TSEV, A.I.: Algebraic systems, Springer, 1973 (translated
from the Russian).
[2] COHN, P.M.: Universal algebra, Reidel, 1981.
[3] NEUMANN, H.: Varieties of groups, Springer, 1967.
[4] BlRKHOFF, G.: Lattice theory, Colloq. Pub!., 25, Amer. Math.
Soc., 1973.
[5] PERKINS, P.: 'Bases of equational theories of semigroups', J. of
Algebra 11, no. 2 (1968), 298-314.
[6] MURSKIi, V.L.: 'The existence in three-valued logic of a closed
class with finite basis, not having a finite system of identities',
Soviet Math. Dokl. 6, no. 4 (1965), 1020-1024. (Dokl. Akad.
Nauk SSSR 163, no. 4 (1965), 815-818)
[7] J6NSSON, B.: 'Algebras whose congruence lattices are distribLl-
tive', Math. Scand 21 (1967), 110-121.
[8] BAKER, K.A.: 'Equational classes of modular lattices', Pacific
J. Math. 28 (1969),9-15.
[9] ScHWABAUER, R: 'A note on commutative semi-groups', Proc.
Amer. Math. Soc. 20 (1969), 503-504.
[10] BAKER, K.A.: 'Primitive satisfaction and equational problems
for lattices and other algebras', Trans. Amer. Math. Soc. 190
(1974), 125-150.
D.M. Smirnov
Editorial comments. A categorical characterization of
varieties of algebraic systems was introduced by F.w.
Lawvere [A1]; for a detailed account of this approach see
[A2].
References
[A1] LAWVERE, F.W.: 'Functional semantics of algebraic theories',
Proc. Nat. Acad. Sci. USA 50 (1963), 869-873.
[A2] MANES, E.G.: Algebraic theories, Springer, 1976.
AMS 1980 Subject Classification: 08BXX
ALGEBRAIC TOPOLOGY The branch of
mathematics in which one studies such properties of
geometrical figures (in a wider sense, of all objects for
which one can speak of continuity), and their mappings
into each other, which remain unchanged under con-
tinuous deformations (homotopies). In principle, the
objective of algebraic topology is a complete listing of
such properties. The very name of algebraic topology
originates from the decisive role of algebraic notions
and algebraic methods in solving problems in this field.
The most important classes of objects whose properties
are studied in algebraic topology include complexes
(polyhedra, cf. Complex), which may be simplicial, cel-
lular, etc.; manifolds (cf. Manifold), which may be
open, closed or with boundary, and may in tum be
divided into smooth (differentiable), analytic, complex-
analytic, piecewise-linear or topological; and fibre bun-
ALGEBRAIC TOPOLOGY
dIes (fibrations, cf. Fibration) and their sections. The
principal types of mappings considered in algebraic
topology are arbitrary continuous, piecewise-linear and
smooth mappings, and their most important subclasses
are: homeomorphisms, in particular, continuous,
piecewise-linear or smooth (diffeomorphisms); imbed-
dings of one object into another and also immersions
(local imbeddings). (cf. Homeomorphism; Diffeomor-
phism.)
A very important notion in algebraic topology is that
of a deformation. The principal types of deformations
include: a homotopy, i.e. an arbitrary continuous
(smooth, piecewise-linear) deformation of continuous
mappings; an isotopy (continuous, smooth, piecewise-
linear), i.e. a deformation of homeomorphisms, imbed-
dings or immersions where in the process of deforma-
tion the mapping is a homomorphism, an imbedding or
an immersion at any moment of time.
The principal internal problems of algebraic topology
include the problem of the classification of manifolds
by homeomorphisms (continuous, smooth, piecewise-
linear), the classification of imbeddings (or immersions)
with respect to isotopies (regular homotopies), and the
classification of general continuous mappings up to
homotopy. An important intermediate role in the solu-
tion of these problems is played by the problem of the
classification of complexes or manifolds by the so-
called homotopy equivalence or homotopy type.
The following, somewhat special, problems played an
important role in the development of algebraic topol-
ogy.
I) The problem of imbedding is usually considered
not in its general form, but rather as imbedding into a
Euclidean space. A very important special case is knot
theory (and the theory of links) in three-dimensional
space, which was one of the origins of algebraic
topology; braid theory is a related case.
2) An important role in the history of algebraic
topology was played by the theory of homology invari-
ants of the position of various sets in a Euclidean space
and the duality laws (cf. Duality in algebraic topology),
connecting the homology of a set and its complement.
3) A number of fundamental results concerned the
computation of the algebraic number of fixed points of
a mapping of a manifold into itself. Many fundamental
facts were discovered for the fixed points of compact
smooth groups of transformations, including cyclic
groups of finite order.
4) Of major technical importance in the development
of algebraic topology were methods developed to solve
the problems on so-called cobordism: Does a manifold
with boundary (cobordism) with a given closed mani-
fold as boundary exist? Problems of this kind originally
arose in the context of the corp.putation of the homo-
123
ALGEBRAIC TOPOLOGY
topy groups of spheres (cf. Homotopy group). Impor-
tant cases of the cobordism problem are solved in the
language of characteristic classes (cf. Characteristic
class).
5) Many facts are now available concerning the
homology invariants of the singularities of vector fields,
frame fields and tensor fields, as well as of the singular-
ities of smooth mappings of manifolds, in particular
into Euclidean space. The solution of this problem
leads, in particular, to characteristic classes. An espe-
cially important case is that of stationary points of
smooth functions on manifolds or of various function-
als on path spaces (extremals); their connection with
homology theory is important in the clarification of the
geometrical structure of manifolds and in obtaining
lower bounds for the number of extremals.
6) The study of the algebraic-topological properties
of an important kind of special spaces - Lie groups -
is closely connected with their algebraic structure, their
representations and the calculus of variations on Lie
groups (cf' Lie group). The results obtained for the
topological structure of Lie groups form the base of
numerous methods and facts of algebraic topology
which apply to arbitrary manifolds. Algebraic topology
of homogeneous manifolds is closely related to methods
of Lie groups.
7) A major role in algebraic topology is played by
special invariants connected with various algebraic
structures over the fundamental group. The simplest
invariants of this type appeared in knot theory and in
the theory of three-dimensional manifolds; the develop-
ment of their algebraic theory subsequently made sub-
stantial progress, and eventually formed a separate dis-
cipline - stable algebra or algebraic K-theory.
8) The analysis of the geometrical structure of a very
large number of examples of the simplest and most fre-
quently encountered manifolds (e.g. Lie groups, homo-
geneous spaces, manifolds of line elements and mani-
folds with discrete groups of motions), in conjunction
with the fundamental principles of Riemannian
geometry, led to the notion of a fibre bundle, which
consists of a space (the total or fibre space), a base
space, a projection of the total space onto the base
space, a fibre (homeomorphic to a 'fibre' of the projec-
tion), and the structure group of fibre transformations.
In this context, one of the central problems in algebraic
topology is the problem of the classification with
respect to homotopy of fibre bundles and their sections.
Principal fibre bundles and vector bundles are espe-
cially important (cf' Principal fibre bundle; Vector bun-
dle).
A universal method for the solution of all principal
problems in algebraic topology involves the construc-
tion of algebraic invariants which are effectively calcul-
124
able in actual examples. and which assume some
discrete set of values; the value of the invariant may
not vary during deformations in the corresponding
class of mappings for the study of which the invariant
has been constructed. A large number of essential
invariants, the abundance of algebraic links between
them and the difficulty in their calculation have deter-
mined the features of modem algebraic topology.
The calculation of algebraic-topological invariants of
the simplest and most frequently encountered mani-
folds is not always a simple matter. Thus, the calcula-
tion of homology invariants of Lie groups and of many
homogeneous spaces required much effort and involved
the use of complicated methods. The calculation of
homotypy groups is even more difficult. A number of
the most important homotypy groups for Lie groups
was calculated, against expectation, using the varia-
tional theory of geodesics; the knowledge of the table
of these homotypy groups for Lie groups made the
classification of certain vector bundles possible.
Most algebraic-topological invariants are a so-called
functor on the category of topological spaces of the
type studied. This means, roughly speaking, that the
values of the invariant undergo natural transformations
when the spaces are mapped into each other. For
example, the fundamental groups of arbitrary spaces or
their homology (cohomology) groups (rings) are con-
nected by homomorphisms induced by continuous
mappings; the characteristic classes (distinguished ele-
ments in (co)homology groups) are mapped into each
other by the homomorphism corresponding to a mor-
phism of manifolds; the result of cohomology opera-
tions performed on an element of the homology (coho-
mology) is converted, after a continuous mapping of
the space, into the result of this operation performed
on the image of this element, etc.
One of the principal properties, on which the study
and the application of almost-all algebraic-topological
invariants are based, is the fact that their effective con-
struction is usually connected with an essential comple-
mentary geometrical structure. For instance, the con-
struction of all the principal invariants up to
homeomorphism for complexes involves a subdivision
into simplices or cells, whereas the result of such a con-
struction has to be invariant with respect to all continu-
ous homeomorphisms, and even with respect to homo-
topy equivalences; examples are the fundamental
group, the Euler characteristic, the homology group
(Betti group), the cohomology ring and the cohomology
operations (cf. Cohomology operation). In a similar
manner, the structure of all the basic homeomorphism
invariants for smooth manifolds involves their prelim-
inary triangulation, i.e. reduction to complexes, or else
to a significant extent the use of analytical tools, such

as the construction of a homology ring by way of dif-
ferential forms (skew-symmetric tensors) and differen-
tial operations on them, or the construction of charac-
teristic classes using the singularities of vector fields,
frame fields or tensor fields. Moreover, it may also be
necessary to use tools of Riemannian geometry in cer-
tain cases, e.g. the definitions of the characteristic
classes of a manifold or of a fibre bundle in terms of
the Riemannian curvature play an important role, even
though the result is invariant wih respect to all continu-
ous homeomorphisms. In this context, the appearance
of basic, effectively calculable invariants in the history
of topology involved the difficult problem of proving
the invariance of these quantities. This is another type
of problem dealt with by algebraic topology. Thus,
rational characteristic classes (or integrals of classes
over cycles) proved to be topologically invariant and
homotopically non-invariant. The totality of charac-
teristic classes over the integers proved to be non-
invariant with respect to continuous or even piecewise-
linear homeomorphisms.
On the contrary, if an invariant is so constructed that
it is invariant (by definition) with respect to continuous
homeomorphisms (or some wider class of transforma-
tions - homotopy equivalences), then the calculation
of such an invariant is usually difficult. The most
important invariants of this type are the homotopy
groups of classes of homotopic mappings of a sphere
into the space under study. The calculation of even the
simplest examples of homotopy groups involves major
difficulties. Thus, the homotopy groups of the spheres
themselves are only incompletely known, despite the
large number of methods which were proposed for such
calculations in the course of development of algebraic
topology (cf. Homotopy group).
The method of classification of classes of homotopic
mappings is based on homology theory in conjunction
with the theory of fibre bundles with its apparatus of
spectral sequences, as well as cohomology operations
and their generalizations. The algebraic discipline which
arose on the basis of the complicated computational
tools of algebraic topology is known as homological
algebra.
All the basic primary constructions of homology
theory for complexes and smooth manifolds by way of
triangulation or differential forms are effectively com-
binatorial - algebraic or analytic. However, it was
noted in the course of development of algebraic topol-
ogy that homology theory can be completely defined by
a small number of formal properties (axioms) on which
its computational tools are based:
1) Homotopic invariance of homology groups.
2) The homology groups are homomorphically
mapped into each other under continuous mappings of
ALGEBRAIC TOPOLOGY
spaces, and products of mappings correspond to the
products of these homomorphisms (functoriality).
3) A certain form of interconnection is postulated
between the cohomology of a complex, a subcomplex
and a quotient complex (exactness axiom).
4) The so-called excision axiom.
5) Normalization: Only the homology groups for a
point need be known, and they must be known to van-
ish in non-zero dimensions.
It was subsequently noted that objects of an alto-
gether different geometrical nature may display all of
the above properties, except for normalization. Such
objects were called generalized or extra-ordinary
homology theories. They were used to improve the cal-
culation methods of algebraic topology. The most
important example is K-theory, which is based on vec-
tor bundles over the space studied instead of their
cycles or differential forms which served for the con-
struction of the ordinary homology groups. Another
important example is cobordism (bordism) theory, in
which only the mappings of closed manifolds into the
space under study are taken instead of arbitrary cycles,
and only mappings of manifolds with boundary of
some classes are taken instead of arbitrary boundaries.
The combination of the methods of algebraic topol-
ogy with a small number of pure geometric facts
applied to the structure of manifolds and their
homeomorphism groups eventually resulted in the solu-
tion (which, from the modem point of view, mayor
may not be considered complete) of classification prob-
lems of manifolds with respect to all kinds of
homeomorphisms (smooth, piecewise-linear, continu-
ous), and in the solution of fundamental problems in
the classification of imbeddings and immersions.
Strangely enough, problems connected with three-
dimensional and four-dimensional manifolds, where the
fundamental problems are still (1977) unsolved,
represent an exception.
Since most of the fundamental internal problems
have now been solved, modem efforts in algebraic
topology tend to concentrate on the application of
these ideas in other fields.
References
[1] ATIYAH, M.F.: K-theory: lectures, Benjamin, 1967.
[2] SEIFERT, H. and TmlELFALL, W.: Lehrbuch der Topologie, Chel-
sea, reprint, 1980.
[3] MILNOR, J.: Morse theory, Princeton Univ. Press, 1963.
[4] MILNOR, J.: Lectures on the h-cobordism theorem, Princeton
Univ. Press, 1965.
[5] MILNOR, J.: Singular points of complex hypersurfaces, Princeton
Univ. Press, 1968.
[6] MILNOR, J. and STASHEFF, J.: 'Characteristic classes', in Ann.
of Math. Studies, Princeton Univ. Press, 1974.
[7] SwITzER, R.M.: Algebraic topology - homotopy and homology,
Springer, 1975.
[8] Hu, S.-T.: Homotopy theory, Academic Press, 1959.
125
ALGEBRAIC TOPOLOGY
[9] STEENROD, N.E.: The topology offibre bundles, Princeton Univ.
Press, 1951.
[10] NOMIZU, K.: Lie groups and differential geometry, Math. Soc.
Japan, 1956.
[11] LERAY, J.: 'Le calcule differentiel et integrale sur une variete
analytique complexe. (Probleme de Cauchy III)" Bull. Soc.
Math. France 87 (1959),81-180.
[12] CHERN, S.S.: Complex manifolds without potential theory,
Springer, 1979.
[13] SPRINGER, G.: Introduction to Riemann surfaces, Addison-
Wesley, 1957.
[14] ALEKsANDROV, P.S.: Combinatorial topology, Graylock, Roches-
ter, 1956 (translated from the Russian).
[15] PONTRYAGIN, L.S.: Grundzuge der kombinatorischen Topologie,
Deutsch. Verlag Wissenschaft., 1956 (translated from the Rus-
sian).
[16] PONTRYAGIN, L.S.: Smooth manifolds and their applications in
homology theory, Moscow, 1976 (in Russian).
[17] SPANIER, E.H.: Algebraic topology, McGraw-Hill, 1966.
[18] NOVIKOV, S.P.: 'Topology', in R.V. Gamkrelidze (ed.): Modem
mathematics,fundamental directions, Vol. I, To appear
(translated from the Russian). (Original: Moscow, 1986).
[19] NOVIKOV, S.P.: 'The Cartan-Serre theorem and intrinsic
homology', Russian Math. Surveys 21, no. 5 (1966), 209-224.
(Uspekhi Mat. Nauk 21, no. 5 (1966
[20] HlRzEBRUCH, F.: TopolOgical methods in algebraic geometry,
Springer, 1978 (translated from the German).
[21] DUBROVIN, B.A., FOMENKO, A.T. and NOVIKOV, S.P.:
'Methods and applications', in Modem geometry, Vol. 1-2,
Springer, 1985 (translated from the Russian).
[22] DUBROVIN, B.A., FOMENKO, A.T. and NOVIKOV, S.P.:
'Methods of homology theory', in Modem geometry, Vol. 3,
Springer, To appear (translated from the Russian).
S.P. Novikov
Editorial comments. The axioms 1) - 5) for a
(co)homology theory are often known as the
Eilenberg - Steenrod axioms (cf. Steenrod - Eilenberg
axioms), [A3]. Some additional good references are [A1J,
[A2] , [A4] - [A6].
References
[A1] ADAMS, J.F.: Stable homotopy and generalized homology,
Univ. of Chicago Press, 1974.
[A2] BOIT, R. and Tu, L.W.: Differential forms in algebraic topol-
ogy, Springer, 1982.
[A3] ElLENBERG, S. and STEENROD, N.: Foundations of algebraic
topology, Princeton Univ. Press, 1952.
[A4] KAROUBI, M.: K-theory. An introduction, Springer, 1978.
[AS] VAISMAN, J.: Cohomology and differential forms, M. Dekker,
1973.
[A6] WHITEHEAD, G.W.: Elements of homotopy theory, Springer,
1978.
AMS 1980 Subject Classification: 55-XX
ALGEBRAIC TORUS - An algebraic group that is
isomorphic over some extension of the ground field to
the direct product of a finite number of multiplicative
groups Gm The group T of all algebraic homomor-
phisms of an algebraic torus T in Gm is known as the
character group of T; it is a free Abelian group of a
rank equal to the dimension of T. If the algebraic torus
T is defined over a field k, then T has a G-module
structure, where G is the Galois gr~oup of the separable
closure of k. The functor T ~ T defines a duality
126
between the category of algebraic tori over k and the
category of Z-free G-modules of finite rank. An alge-
braic torus over k that is isomorphic to a product of
groups Gm over its ground field k is called split over k;
any algebraic torus over k splits over a finite separable
extension of k. The role played by algebraic tori in the
theory of algebraic groups greatly resembles the role
played by tori in the theory of Lie groups. The study of
algebraic tori defined over algebraic number fields and
other fields, such as finite fields, occupies an important
place in problems of arithmetic and in the classification
of algebraic groups. Cf. Linear algebraic group;
Tamagawa number.
References
[1] BOREL, A.: Linear algebraic groups, Benjamin, 1969.
[2] ONO, T.: 'Arithmetic of algebraic tori', Ann. of Math. (2) 74,
no. 1 (1961), 101-139.
[3] ONO, T.: 'On the Tamagawa number of algebraic tori', Ann. of
Math. (2) 78, no. 1 (1963),47-73. VE T7 k k'-
y os resens II
AMS 1980 Subject Classification: 20G15
ALGEBRAIC VARIETIES, ARITHMETIC OF, arith-
metical algebraic geometry - The branch of algebraic
geometry in which one studies properties of algebraic
varieties defined over fields of so-called arithmetic type,
i.e. finite, local and global fields of algebraic numbers
or algebraic functions. In the case of a finite field, its
main subject is the study of the number of rational
points of the algebraic variety in this field and in finite
extensions of it. The zeta-function of the variety which
is used in these studies strongly influenced the develop-
ment of the methods of algebraic geometry. Estimates
from below of the number of (rational) points [1], [4]
are also important.
If X is an algebraic variety (or scheme) over a local
field K with field of residues k, then the study of the
set X (K) of rational points with values in K combines
two different problems: To find solutions of
congruences (or points of the variety over a finite field)
and to find integral or rational solutions of Diophan-
tine equations (cf. Hasse principle). If the variety X is
defined by a set of equations with coefficients from the
ring A of integral elements of the field K, then it is pos-
sible to define the reduction of this variety by the same
set of equations, but with coefficients taken modulo the
maximal ideal of A. A 'variety' X 0 over the field of
residues k and a canonical mapping, or reduction,
Red: X(K) __ Xo(k)
are obtained. This description of the reduction is diffi-
cult to explain in terms of classical algebraic geometry.
This was one of the reasons for the introduction of the
concept of a scheme - a language suitable for a
rigorous description of the process. The main problem

is to determine the iniage of the mapping Red, i.e. to
find the points x EXo(k) that come from the rational
K-points of the variety. The Hensel lemma states that x
is such a point if x is a non-singular point. For more
general results on this subject see [4].
Another type of problems concerned with the local
arithmetic of algebraic varieties is the study of forms
over such fields. Let F be a form in n variables of
degree d over a local field; Artin's conjecture is that if
n>d2 , then the equation F=O has a non-trivial solu-
tion. It is known that this statement is true in the case
of function fields. It was proved for p-adic fields that
for each d there exists a finite number A (d) of primes
such that Artin's conjecture holds for forms of degree
d, if P gA(d). It was shown in 1966 that the set A (4) is
non-empty, which proved Artin's conjecture to be false
[4]. It is not yet (1977) known whether or not it is valid
for forms of odd degree.
The arithmetic of algebraic varieties over global
fields is the largest and most diverse part of algebraic
geometry. It includes Diophantine geometry, class field
theory, the theory of zeta-functions of varieties, and
complex multiplication of Abelian functions (or
varieties). All these theories were developed in parallel
for number and function fields. Such a possibility was
first demonstrated by the development of class field
theory in the nineteen-thirties; it is based on the rough
analogy between these fields which is most clearly man-
ifested in the construction of the theory of schemes.
References
[1) BOREVICH, Z.I. and SHAFAREVICH, I.R.: Number theory, Acad.
Press, 1966 (translated from the Russian).
[2) WElL, A.: 'Number theory and algebraic geometry', in Proc.
Internat. Math. Congress Cambridge, 1950, Vol. 2, 1950, pp.
90-100.
[3) GROTHENDIECK, A. and DIEUDONNE, 1.: 'Elements de
geometrie algebrique', Publ. Math. IHES 4;8;11;17;20;24;28;32.
[4) PARSHIN, A.N.: 'Arithmetic on algebraic varieties', J. Soviet
Math. 1, no. 5 (1973), 594-620. (Itogi Nauk. Algebra Topol.
Geom 19701 (1970), 111-152)
[5) SWINNERTON-DYER, H.P.F.: 'Applications of algebraic
geometry to number theory', in Proc. 1969 summer inst. number
theory, Proc. of Symp. Pure Math., Vol. 20, Amer. Math. Soc.,
1971.
A.N. Parshm
. geometry (resolution of singularities; the moduli prob-
AMS 1980 Subject Classification: 14J20, 14K15, 14GXX
ALGEBRAIC VARIETY - One of the principal
objects of study in algebraic geometry. The modern
definition of an algebraic variety as a reduced scheme
of finite type over a field k is the result of a long evolu-
tion. The classical definition of an algebraic variety was
limited to affine and projective algebraic sets over the
fields of real or complex numbers (d. Affine algebraic
set; Projective algebraic set). As a result of the studies
initiated in the late nineteen-twenties by B.L. van der
Waerden, E. Noether and others, the concept of an
ALGEBRAIC VARIETY
algebraic variety was subjected to significant algebrai-
zation, which made it possible to consider algebraic
varieties over arbitrary fields. A. Wei! [6] applied the
idea of the construction of differentiable manifolds by
glueing to algebraic varieties. An abstract algebraic
variety is obtained in this way and is defined as a sys-
tem (Va) of affine algebraic sets over a field k, in each
one of which open subsets WaP C Va, corresponding to
the isomorphic open subsets W,Ba C V,B, are chosen. All
basic concepts of classical algebraic geometry could be
transferred to such varieties. Examples of abstract alge-
braic varieties, non-isomorphic to algebraic subsets of a
projective space, were subsequently constructed by M.
Nagata and H. Hironaka [2], [3]. They used complete
algebraic varieties (d. Complete algebraic variety) as the
analogues of projective algebraic sets.
J.-P. Serre [5] has noted that the unified definition of
differentiable manifolds and analytic spaces as ringed
topological spaces has its analogue in algebraic
geometry as well. Accordingly, algebraic varieties were
defined as ringed spaces (cf. Ringed space), locally iso-
morphic to an affine algebraic set over a field k with
the Zariski topology and with a sheaf of germs of regu-
lar functions on it. The supplementary structure of a
ringed space on an algebraic variety makes it possible
to simplify various constructions with abstract algebraic
varieties, and study them using methods of homological
algebra which involve sheaf theory.
At the International Mathematical Congress in Edin-
burgh in 1958, A. Grothendieck outlined the possibili-
ties of a further generalization of the concept of an
algebraic variety by relating it to the theory of schemes.
After the foundations of this theory had been esta-
blished [4] a new meaning was imparted to algebraic
varieties - viz. that of reduced schemes of finite type
over a field k, such affine (or projective) schemes
became known as affine (or projective) varieties (cf.
Scheme; Reduced scheme). The inclusion of algebraic
varieties in the broader framework of schemes also
proved useful in a number of problems in algebraic
lem, etc.).
Another generalization of the concept of an algebraic
variety is related to the concept of an algebraic space.
Any algebraic variety over the field of complex
numbers has the structure of a complex analytic space,
which makes it possible to use topological and tran-
scendental methods in its study (d. Kahler manifold).
Many problems in number theory (the theory of
congruences, Diophantine equations, modular forms,
etc.) involve the study of algebraic varieties over finite
fields and over algebraic number fields (d. Algebraic
varieties, arithmetic of; Diophantine geometry; Zeta-
function in algebraic geometry).
127
ALGEBRAIC VARIETY
References
[1) BALDASSARRI, M.: Algebraic varieties, Springer, 1956.
(2) SHAFAREVICH, I.R.: Basic algebraic geometry, Springer, 1977
(translated from the Russian).
(3) DOLGACHEV, I.V.: 'Abstract algebraic geometry', J. Soviet
Math. 2, no. 3 (1974), 264-303. (Itogi Nauk. i Tekhn. Algebra
Topol. Geom 10 (1972),47-112)
(4) GROTIJENDIECK, A. and DIEUDONNE, J.: 'Elements de
geometrie algebrique', Publ. Math. IHES 4 (1960).
(5) SERRE, J.-P.: 'Faiseaux algebriques coherentes', Ann. of Math.
(2) 61, no. 2 (1955), 197-278.
(6) WElL, A.: Foundations of algebraic geometry, Amer. Math. Soc.,
1946.
I. V. Doigachev
AMS 1980 Subject Classification: 14A10
ALGEBRAIC VARIETY, AUTOMORPIDSM OF AN
- An invertible morphism of an algebraic variety (or
scheme) into itself. The group of all automorphisms of
an algebraic variety X, which is usually denoted by
Aut X, is an important invariant of X. Studies of the
action of the automorphism group of an algebraic
variety on objects functorially connected with X such as
the Picard group; the Chow ring; the K-functor, and the
cohomology group are a tool used in the study of the
varieties themselves. The automorphism group of an
algebraic variety is important for the concept of forms
(cf. Fonn) of an algebraic variety. For complete alge-
braic varieties over the field of complex numbers, the
automorphism group is identical with the group of
biholomorphic automorphisms.
The structure of the group Aut X is known for a
number of simple algebraic varieties. For instance, if X
is projective n-dimensional space pn over a field k,
then anyone of its automorphisms is a linear projective
transformation and Aut pn becomes identical with the
projective linear group PLG(n + 1, k). The automor-
phism group of an elliptic curve and, in general, of any
Abelian variety A, is an extension of the group G of
automorphisms which preserve the structure of the
Abelian variety, by the group A (k) of translations in
the points of A, i.e. the sequence of groups
1 ~ A (k) ~ AutA ~ G ~ 1
is exact. If X is a smooth, complete algebraic curve of
genus g> 1, then the group Aut X is finite; an estimate
of its order as a function of g is known (cf. Algebraic
curve). For automorphisms of surfaces, see Algebraic
surface.
In the case of algebraic varieties with an ample
canonical or anti-canonical invertible sheaf the auto-
morphism group is an algebraic subgroup of the group
PLG(N, k) for some N. The automorphism group of a
smooth hypersurface of dimension n;;;' 2: and degree
d~3 is finite [1].
In the above examples, Aut X has the natural struc-
ture of an algebraic group, perhaps with an infinite
128
number of connected components; this is also true in
the general case [2].
Families of automorphisms are considered in the
modern approach to automorphism groups of algebraic
varieties. A family of automorphisms of a variety X
with parameter scheme T is a set of automorphisms of
the product XX T that commute with projection onto
the second factor; the set of families of automorphisms
with parameter scheme T is denoted by AutT(XX T).
One thus obtains a contravariant functor
TI-+AutT(XX T). If the variety X is complete, then this
functor is locally representable (cf. Representable
functor) by an algebraic group scheme with at most a
countable number of connected components [3]. A.
Grothendieck gave a proof of this fact for projective
varieties, and this theorem has been extended to the
case of proper flat schemes of morphisms. The scheme
representing this functor is not necessarily reduced even
if X is a smooth projective surface; however, if the
characteristic of the ground field is zero, or if X is a
smooth curve or a smooth hypersurface, then the con-
nected component of the unit of this scheme is a
variety.
For incomplete varieties the automorphism functor is
not always representable in the category of schemes.
For an affine variety, the automorphism functor is
representable in the category of inductive limits of
schemes.
Apart from the simple case of the affine straight line,
for the affine spaces only the automorphism group of
the affine plane is known. It is a free product of two of
its subgroups with as amalgamated subgroup their
intersection, viz. the subgroup of linear affine transfor-
mations and the subgroup of triangular automorphisms,
i.e. transformations of the form
x' = ax+b,
y' = cy +f(x},
where a, b, cEk, a*O, c*O, while f(x) is an arbitrary
polynomial in x [4], [5]. For the treatment of affine
algebraic surfaces transitively acted upon by the auto-
morphism group, see [6].
References
[II MATSUMURA, H. and MONSKY, P.: 'On the automorphisms of
hypersurfaces', J. Math. Kyoto Univ. 3 (1964),347-361.
[2) MATSUSAKA, T.: 'Polarized varieties, fields of moduli and gen-
eralized Kummer varieties of polarized Abelian varieties',
Amer. J. Math. 80 (1958), 45-82.
[31 MATSUMURA, H. and OORT, F.: 'Representability of group
functors and automorphisms of algebraic schemes', Invent.
Math. 4 (1967), 1-25.
[4] ENGEL, W.: 'Ganze Cremona-Transformationen von Prim-
zahlgrad in der Ebene', Math. Ann. 136 (1958),319-325.
[5] SHAFAREVICH, I.R.: 'On some infinite-dimensional groups',
Rend. di Mat. e Appl. 25 (1966),208-212.
[6] GIZATULLlN, M.KH.: 'Quasi homogeneous affine surfaces',

Math. USSR-Izv. 5, no. 5 (1971), 1057-1082. (lzv. Akad. Nauk
SSSR Ser. Mat. 35 (1971), 1047-1071)
[7] ROTH, L.: Algebraic threefolds, Springer, 1955.
VI. Danilov
VA. Iskhovskikh
AMS 1980 Subject Classification: 14JXX, 14A 10
ALGEBRAICALLY CLOSED FIELD - A field k in
which any polynomial of non-zero degree over k has at
least one root. In fact, it follows that for an algebrai-
cally closed field k each polynomial of degree n over k
has exactly n roots in k, i.e. each irreducible polynomial
from the ring of polynomials k[x] is of degree one. A
field k is algebraically closed if and only if it has no
proper algebraic extension (d. Extension of a field). For
any field k, there exists a unique (up to isomorphism)
algebraic extension of k that is algebraically closed; it
is called th~ algebraic closure of k and is usually
denoted by k. Any algebraically clo~ed field containing
k contains a subfield isomorphic to k.
The field of complex numbers is the algebraic closure
of the field of real numbers. This is the fundamental
theorem of algebra (d. Algebra, fundamental theorem
of).
References
[I] O. and SAMUEL, P.: Commutative algebra, I, Springer,
ZARISKI,
1975.
[2] LANG, S.: Algebra, Addison-Wesley, 1974.
o.A. Ivanova
AMS 1980 Subject Classification: 12F05
ALGOL - A general name for a number of algo-
rithmic languages (d. Algorithmic language) used for
automatic programming and for the publication of algo-
rithms (it is an abbreviation of 'ALGOrithmic
Language').
The first variant of Algol was developed in 1958 by
an international scientific team. The 1960 International
Conference in Paris adopted the language 'Algol-60',
which combined several useful features of the program-
ming languages then known. Algol-60, and its descen-
dent like Pascal, gained widespread use and the name
Algol now usually refers to 'Algol-60'. It is especially
suited for the description of algorithms of numerical
analysis. The language is machine-independent and
does not provide standard ways of performing input or
output operations. Special concrete versions of standard
Algol may be developed for individual computers, and
each one is a language accepted by the compiler of the
individual machine. As a rule, the transition from the
standard language to its concrete version is natural and
does not involve much labour. The basic symbols in
Algol are decimal numbers, the upper-case and lower-
case letters of the English alphabet, punctuation marks,
ALGOL-68
symbols of arithmetical and logical operations, other
special symbols and certain English words (e.g., 'begin',
'end', 'real', 'integer', 'array'). From the basic symbols
of the language further entities are formed according to
certain rules, such as numbers, identifiers (names), sim-
ple variables, subscripted variables, function designa-
tors, expressions, statements, and comments. There are
several basic types of statements: an assignment state-
ment, a goto statement, a conditional statement which,
depending on the outcome of the logical expression
contained in it, actuates one of its own inner state-
ments, and a for statement (a loop). A group of state-
ments may be combined into a compound statement or
a block (which contains declarations). The Algol nota-
tion of an algorithm may include definitiOns of pro-
cedures. The definition of a procedure includes a head-
ing and a body. The body of a procedure may be a
statement (most often a block) described in the usual
Algol notation. In order to render the Algol language
more flexible, procedures expressed in some other
language (e.g. machine language) may be employed. A
procedure may be called with the aid of a procedure
statement consisting of the procedure identifier and a
list of actual parameters corresponding to the formal
parameters provided in the heading of the procedure
with the given identifier. An expression may contain
function designators, which call the procedure for the
computation of a value. The procedure call may be
recursive, i.e. a call which provides for the re-call of the
same procedure while executing it. The language poten-
tialities of standard Algol are often restricted in its con-
crete implementations. The proposed successor to
Algol-60 is the language A1gol-68, which has a signifi-
cantly different structure, contains many new concepts
and possibilities, and is designed for more powerful
machines.
References
[lA] NAUR, P., ET AL.: 'Revised report on the algorithmic language
Algol 60', Num. Math. 4 (1963), 420-453.
[IB] NAUR, P., ET AL.: 'Revised report on the algorithmic language
Algol 60', Comm ACM 6 (1963), 1-17.
[lC] NAUR, P., ET AL.: 'Revised report on the algorithmic language
Algol 60', The Computer Journal 5 (1963), 349-367.
[2] LAVROV, S.S.: A universal programming language (Algol-60),
Moscow, 1967 (in Russian).
[3A] WUNGAARDEN, A. VAN, ET AL.: 'Report on the algorithmic
language Algol 68', Num. Math. 14 (1969), 79-218.
[3B] WUNGAARDEN, A. VAN, ET AL.: 'Revised report on the algo-
rithmic language Algol 68', Acta Inform. 5 (1975), 1-236.
V V Martynyuk
AMS 1980 Subject Classification: 68BXX
ALGOL-68 - A universal algorithmic language. It
was developed during 1964 - 1968 by a team of scien-
tists from 12 countries, in the framework of the work-
ing group on Algol of the International Federation for
129
ALGOL-68
Information Processing, for the exchange of algorithms,
their effective realization by various computers, and as
a means of their study. While similar in style to Algol-
60, Algol-68 differs substantially from it in having
many more constructions, which are also of a more
general nature. The basic data types, in addition to the
Algol-60 'real', 'integer' and 'Boolean', are 'character'
(for alpha-numeric information), 'format' (to describe
the format of the external data), names and routines
(procedures). Thus, names and procedures can be 'com-
puted' when a program is executed in Algol-68, even
though such a computation is limited by the necessity
of selecting the value of, for example, a procedure out
of a given finite population of 'routine texts'. The basic
types may be used to construct new, composite types
by induction; these types may be homogeneous index-
able sequences of data of the same type (multiple
values), or else ordered lists of data of arbitrary type
(structured values).
In addition to the conventional approach for defining
procedures, Algol-68 contains means for defining so-
called infix operators of the type x +y. The presence of
a priority definition makes it possible to establish prior-
ity relations between the infix operators introduced.
The identity definition, which is typical of Algol-68, is a
general construction for defining variables, for assign-
ing initial values, for the transmission of actual parame-
ters in procedures and for creating synonyms.
An expression in Algol-68 may contain an assign-
ment statement, or any sequence of statements that
produces a value. Together with the possibility of com-
puting names and procedures, and also by the intro-
duction of parenthesis pairs for conditional expressions,
this permits Algol-68 constructions such as illustrated
by the following example:
1) Algol-68:
if x >0 then u else z Ii
: = a +(m<n Isin Icos)(t : = xj2),
2) Algol-60:
t := xj2;
r : = a + if m <n then sin(t) else cos(t);
if x>O then u: =r else z: =r.
An Algol-68 program consists of closed, serial, condi-
tional and collateral clauses. The first three types of
clauses generalize such Algol-60 concepts as block,
compound statement and conditional expression and
statement. Collateral clauses signify non-ordered sets of
component phrases and are used, in particular, to indi-
cate parallel branches in the course of execution of the
program.
The description of the semantics of Algol-68 is
characterized by a profound analysis of the principal
concepts of algorithmic languages, which makes it pos-
130
sible to give a description of the execution of a pro-
gram using of a small number of independent funda-
mental concepts. The objects may be external (with
respect to the structure of the program) and internal
(data, including procedures and names). The relations
between external (E) and internal (I) objects are intro-
duced as axioms, e.g. 'E) contains ~', 'E) identifies
~', 'E has 1', 'I) names 12 ', ' I) is a component of 12 ',
etc. The execution of a program is described in terms
of the relations thus introduced as a function of the
analysis of the program.
A typical syntactic feature of Algol-68 is that it is
given in the form of a two-level grammar, in which the
Algol-68 production rules are themselves well-formed
texts in some meta-language, defined by its own gen-
erating grammar. For example, grammatical rules of
Algol-68 may have the following form:
chain of NOTIONs separated by SEPARATORs:
NOTION;
NOTION, SEPARATOR,
chain of NOTIONs separated by SEPARATORs.
reference to MODE assignation:
reference to MODE destination,
becomes symbol, MODE source.
The capitalized words, called 'metanotions', are the
grammatical units of the meta-language, whose produc-
tion rules may have forms such as:
NOTION: identifier; operator.
SEPARATOR: comma; semi-colon.
MODE: integral; Boolean; reference to MODE.
Some of the concepts of the meta-language, such as
MODE, can have an infinite number of productions.
The proper production rules of Algol-68 are obtained
by systematically replacing the metanotions of the
meta-language in the syntax rules by some terminal
production of these metanotions. The resulting rules, in
the meta-linguistic notation of Algol-60, may have, for
example, the following form:
< chain of identifiers separated by commas > :: =
< identifier > I < identifier >, < chain of identifiers
separated by commas >
< reference to integer assignation> :: = < refer-
ence to integer destination > : = < integer source >.
The use of a two-level grammar makes it possible,
first, to reduce the number of production rules of the
same type and, secondly, to give a syntactic expression
to the attributive information of the notions and to
ensure contextual relations which must otherwise be
formulated as context conditions (static semantics).
References
[lA) WUNGAARDEN, A. VAN, ET AL.: 'Report on the algorithmic
language Algol 68', Num. Math. 14 (1969), 79-218.
[I B) WUNGAARDEN, A. VAN, ET AL.: 'Revised report on the algo-
rithmic language Algol 68', Acta Inform. 5 (1975), 1-236.

[2) LINDSEY, C.H. and MEULEN, S.G. VAN DER: Informal introduc-
tion to Algol 68, North-Holland, 1977.
A.P. Ershov
AMS 1980 Subject Classification: 68BXX
ALGORITHM - Detailed instructions defining a
computational process (which is then said to be algo-
rithmic), which begins with an arbitrary input (out of a
certain number of inputs which are possible for the
given algorithm), and with instructions aimed at
obtaining a result (or output) which is fully determined
by the input. For instance, the rules taught in elemen-
tary schools for column-wise addition, subtraction, mul-
tiplication and division are algorithms; in these algo-
rithms the possible results are non-negative integers
written in the decimal system, while the possible inputs
are ordered pairs of such numbers. It is, in general, not
assumed that the result is necessarily obtained: the pro-
cess of applying an algorithm to some possible input
(i.e. an algorithmic process which proceeds starting
from this input) may also terminate without a result
being obtained (in such a case one has a no-result stop)
or may not terminate at all. If the process terminates
(does not terminate), then one says that the input is
suitable (is not suitable) for the algorithm in question.
An important result in this area is the undecidability
of the so-called halting problem. It is possible to con-
struct an algorithm a such that there is no algorithm to
decide whether or not some arbitrary input which is
possible for a is also suitable for a. In particular, one
can find such an algorithm a for which the set of its
possible inputs is the set of natural numbers.
The notion of an algorithm is one of the central con-
cepts in modem mathematics, especially in the area of
computations. An example is given below. Finding
the numerical solution of equations of a given type is
tantamount to constructing an algorithm which con-
verts an arbitrary equation of this type and an arbitrary
positive rational number t: into a number (or a set of
numbers) which differs from the root(s) of this equation
by less than t:. However, the term 'computational pro-
cess', as used in the context of an algorithm, must not
be understood to mean merely numerical calculations:
even in elementary algebra computations with letters
are used. Standard arithmetic calculations involve some
symbols which are not numbers (parentheses, equality
signs, symbols for arithmetic operators). It is accord-
ingly expedient to consider algorithms involving arbi-
trary symbols and objects composed from symbols. The
simplest example of such an object is a linear sequence
of symbols forming a word. It is also possible to con-
sider 'non-linear' objects - algebraic matrices, deriva-
tion trees of some formal language and general graphs.
Intuitively, the least requirement on inputs and results
of algorithms is that they must be constructive objects
ALGORITHM
(d. Constructive object). Thus, the concept of an algo-
rithm is very general. One may speak of an algorithm
for the translation from one language into another, of
an algorithm of an air traffic controller (processing
information on the motions of aircraft according to
fixed rules), and other examples of algorithmically
described control processes. Consequently, the concept
of an algorithm is one of the central ideas in cybernet-
ics and computer science.
Example of an algorithm. Let the possible inputs and
the possible results be all possible words over the alpha-
bet {a, b }. The transition from a word X to a word Y
is 'permissible' in the following two cases (where P
denotes an arbitrary word): 1) X is of the form aP, and
Y is of the form Pb; or 2) X is of the form baP, while
Y is of the form Paba. The instructions are formulated
as follows: 'designating some word as the input, exe-
cute the permitted transitions until a word of the form
aaP is obtained; then halt and let word P be the
result'. These instructions constitute an algorithm,
which will be denoted by m. Take the word babaa as
the input. After one transition one obtains baaaba;
after the second, one obtains aabaaba. Here the algo-
rithm halts with baaba as result. Consider another
input baaba. One obtains, in succession,
abaaba, baabab, abababa, bababab,babababa,.... It
can be shown that this process will never end (i.e. none
of the generated words will ever begin with aa). Now
take abaab as input. One obtains
baabb, abbaba, bbabab; no further permissible transi-
tions are possible, and, at the same time, the termina-
tion condition is not satisfied. One has a no-result stop.
Thus, input babaa is suitable for m, and inputs baaba
and abaab are not.
The significance of algorithms. In science, algorithms
are encountered everywhere: a 'general' solution to a
problem requires an algorithm. The ability of man to
add numbers is an example of this. Not the fact that he
can find, sooner or later, the sum of any two numbers,
but rather in the sense that he possesses a method for
the unique addition of any two numbers given in writ-
ten form, i.e. an addition algorithm such as the well-
known column-wise addition method. 'Generally' solv-
ing a problem is formalized as solving an algorithmic
problem. An algorithmic problem consists of separate
problem instances, and a single algorithm needs to be
found for the solution of all of them. If there is no
such algorithm, one says that the algorithmic problem
is unsolvable. Thus, the problem of numerically solving
equations of a given type and the problem of automatic
translation are algorithmic problems. Their problem
instances are, respectively, to find the numerical solu-
tion of individual equations of the given type, and the
131
ALGORITHM
translation of individual phrases. Other examples are
the verification of algebraic equalities in algebra; the
algorithmic problem of recognition of the deducibility
of propositions from given axioms in mathematical
logic. (In mathematical logic the concept of an algo-
rithm is also important since it serves as the basis of
the key concept of a calculus, which is a generalization
and a more precise form of the intuitive concepts of
'proof and 'demonstration'.) The establishment of the
unsolvability of a given algorithmic problem (e.g. the
problem of establishing the truth or demonstrability of
sentences in some logical-mathematical language) is
important, because it shows that, in principle, specific
problem instances of that problem can only be solved
by methods which are specific to each individual prob-
lem instance.
The scientific importance of the concept of an algo-
rithm has been recognized for a long time. Since the
earliest times, people have looked for constructive
methods to solve mathematical problems. Such efforts
usually benefitted from the new availability of suitable
symbolic notations. This process has made a signifi-
cant contribution to scientific knowledge, especially so
after it had become clear that some problems are
inherently unsolvable (squaring the circle, etc.). After it
had been recognized that certain problems cannot be
solved by direct calculations, the theoretical concept of
a set was born in the nineteenth century. It is only after
a period of rapid development of this concept (which
did not involve problems of constructive methods in
their modem sense), that it was possible in the mid-
twentieth century to return to constructive problems,
but this was done on a different level, enriched by the
concept of an algorithm which had crystallized in the
meantime. This concept forms the base of the con-
structive trend in mathematics.
The word 'algorithm' comes from the word 'algor-
itmi', the latter being the Latin transliteration of the
name of the nineth century Arab mathematician AI-
Khwarizrni. An 'algorithm' in the Europe of the Middle
Ages was 'the decimal-positional system and the art of
calculating with it', since it is through the Latin transla-
tion (twelfth century) of the treatise of AI-Khwarizmi
that the positional system was introduced in Europe.
The structure of an algorithmic process. An algorithmic
process is a process of consecutive conversion of con-
structive objects by discrete 'steps', each step consisting
of the replacement of one constructive object by
another one. Thus, when applying the algorithm 2{ to
the word baaba, there follow, in succession, the words
baaba, abaaba, baabab, etc. If one applies. say. the algo-
rithm of column-wise subtraction to the pair of
numbers <307, 49>, then one obtains, successively.
the following constructive objects:
132
307 307 307 307
- 49 - 49 - 49 - 49
- - 8 ----s8 258'
It will be noted that, in a series of successive con-
structive objects, each successive constructive object is
fully determined (in the framework of the given algo-
rithm) by the constructive object immediately preceding
it. In the more rigorous approach, it is also assumed
that the transition from any constructive object to the
one immediately following it is sufficiently 'elementary'
- in the sense that the one-step transformation of the
preceding into the immediately succeeding object is
local (it is not the entire constructive object which is
transformed, but only its algorithm-limited part; more-
over, this transformation itself is not determined by the
complete preceding object, but only by a limited part
of it).
Accordingly, one has not only a set of possible
inputs and a set of possible results, but also a set of
possible intermediate results, representing the working
medium in which the algorithmic process is evolving.
Regarding m, all these three sets coincide, but this is
not true of the algorithm of column-wise subtraction -
the possible inputs are pairs of numbers, the possible
results are numbers (all in the decimal system), while
the possible intermediate results are 'three-storey'
records of the type
p
-q
r
where q is the written record of the number in the
decimal system, r is a similar record or the empty word,
while p is the written record of a number in the decimal
system, with dots over some of the digits. As a rule,
several (not independent) parameters which are charac-
teristic of an algorithm may be distinguished: 1) the set
of possible inputs; 2) the set of possible results; 3) the
set of possible intermediate results; 4) the starting rule;
5) the direct transition rule; 6) the termination rule;
and 7) the rule for the retrieval of the result.
Formalizing the concept of an algorithm. The concept
of an algorithm in its general form is a fundamental
mathematical concept which cannot be defined in terms
of simpler concepts. Strictly speaking, formalizations of
the concept of an algorithm actually restrict it some-
what. In each such formalization an exact definition is
given of some class within which each one of the
parameters may vary. The various formalizations differ
from each other by the selection of such classes. Since
the parameters unambiguously define some algorithm. a
choice of parameter domains determines some class of
algorithms. However, such a choice may claim to be a
formalization of the intuitive concept of algorithm. only
if one can be sure that for any 'intuitive' algorithm

there exists an equivalent algorithm in the class of algo-
rithms defined by the choice under consideration. This
requirement is formulated for each formalization as a
fundamental hypothesis, which, in the present state of
knowledge, cannot be mathematically demonstrated.
The first formalizations of this type are due to E.L.
Post [5] and to AM. Turing [3], [4], whose construc-
tions anticipated in many respects the ideas on which
modern computers are based. There are also versions
formulated by AA Markov [10], [11] (cf. Nonnal algo-
rithm) and by AN. Kolmogorov [12], [13], who based
the approach on constructive topological complexes of
a certain type, which makes it possible to express more
precisely the property of 'being local' of a transforma-
tion. In each of these formalizations the fundamental
hypothesis is in satisfactory agreement with practice.
Another argument in favour of this hypothesis is the
fact that, as can be demonstrated, all the formal ver-
sions (of effective computability) which have been pro-
posed are equivalent.
An an example, consider the formalization proposed
by Turing. In its original form, this was a description
of some theoretical computer consisting of 1) an infin-
ite tape, divided into consecutive cells, each cell capable
of containing some symbol out of the 'tape alphabet' of
the machine; and 2) a 'finite control,' which at each
moment of time is in some 'state' (out of a given finite
list of states). The finite control can move along the
tape, one cell at a time, and change the contents of the
cells it passes. A program which monitors the motion
of the finite control constitutes the algorithm for the
computations conducted on such a machine (a 'Turing
algorithm'). For a more exact and more detailed
description see Turing machine. Here a modernized
exposition of Turing's construction is given. To define
a Turing algorithm, one specifies a) pairwise disjoint
alphabets B, D, C, with a distinguished letter A in D
and distinguished letters IX and '" in C; and b) a set of
pairs of the form <p~, T/Tq> where p, qEC,
~,T/EB U D, while T is one of the three symbols
- , 0, +; in this set (which is called a program) there
are no two pairs with the same first coordinate. The
possible inputs and the possible outputs are words over
B, while the possible intermediate results are words
over BUD U C which contain not more than one
letter of C. The starting rule is that the initial word P is
converted to the word AaPA. The termination rule is
that an intermediate result which contains '" is final.
The rule for the retrieval of the output is that the out-
put is the sequence of letters in the final intermediate
result which follows '" and precedes the first letter not
belonging to B. The transformation rule which converts
A into A' is as follows. The letter A is written to the left
and to the right of A; in the word thus formed, the part
ALGORITHM, COMPLEXITY OF DESCRIPTION OF AN
of the form qJ~, where p EC, ~EB U D, is replaced by
a word Q in accordance with the rule: a pair with the
first term p~ is sought in the program; let the second
term of this pair be T/Tq ; if Tis -, Q =qfT/; if T=O,
Q =qT/; if T is +, Q=fT/q. The word produced by this
exchange is A'. For references, see [3], [4], [5], [10],
[11], [12], [13] of Algorithms, theory of.
V.A. Uspenskii
Editorial comments. The fundamental hypothesis men-
tioned above is usually called the Church- Turing thesis or
Turing's hypothesis. This thesis states that what can be
effectively computed according to one's intuition about what
constitutes an effective process, can also be computed by a
Turing machine, i.e. a formally defined notion of algorithm.
By its very nature of identifying an intuitive informal notion
with a formal precise definition, the Church - Turing thesis
cannot be formally proved.
Many formalizations have turned out to be equivalent, and
no formalization has turned out to be strictly more powerful.
Apart from the formalizations mentioned above, the follow-
ing references are the most standard ones in the West:
Church's formalization of numbers and computable func-
tions by lambda terms, Kleene's general recursive schemes,
and the register machine model proposed by J.C.
Shepherdson and H.E. Sturgis, which was extended into a
more realistic computer model by SA Cook and RA Reck-
how. Markov's and Kolmogorov's versions are less well
known. For completeness the original references to Turing
and Post are also given.
References
[A1A] TUluNG, A.M.: 'On computable numbers, with an application
to the Entscheidungsproblem', Proc. London Math. Soc. Ser.
242 (1936), 230-265.
[A1B] TuIuNG, A.M.: 'Corrections to 'On computable numbers,
with an application to the Entscheidungsproblem", Proc.
London Math. Soc. Ser. 243 (1937), 544-546.
[A2] CHuRCH, A.: 'An unsolvable problem of elementary number
theory', Amer. J. Math. 58 (1936), 345-363.
[A3] KLEENE, S.c.: 'General recursive functions of natural
numbers', Math. Ann. 112 (1936),727-742.
[A4] POST, E.L.: 'Formal reductions of the general combinatorial
decision problem', Amer. J. Math. 65 (1943), 197-268.
[A5] SHEPHERDSON, J.C. and STURGIS, H.E.: 'Computability of
recursive functions', J. Assoc. Compo Mach. 10 (1963),
217-255.
[A6] COOK, SA and RECKHOW, R.A.: 'Time-bounded random
access machines', J. Computer and System Sciences 7
(1973),354-375.
AMS 1980 Subject Classification: 03D03, 03D10,
68CXX, 68C05
ALGORITHM, COMPLEXITY OF DESCRIPTION
OF AN, program-size compleXity - A measure which
characterizes the length of description of an algorithm.
The complexity of description of an algorithm is
defined differently, depending on the specific defini-
tion. At present (1987) there is as yet no generally
valid definition of the concept, and the most frequently
occurring cases are reviewed below.
133
ALGORITHM, COMPLEXITY OF DESCRIPTION OF AN
The complexity of description of a nonnal algorithm
is usually understood to mean the length of its encod-
ing, i.e. the length of all its substitution formulas when
written in one line (a special separating character is
interposed between the individual formulas). The com-
plexity of description of a Turing machine usually
means the number of its internal states and external
symbols. A Turing machine may also be characterized
by the number of commands of the machine in ques-
tion. In the case of recursive functions (cf. Recursive
function) given by recursion schemes, the number of
letters in these schemes usually serves as the measure of
their complexity.
An axiomatic definition of the complexity of descrip-
tion of an algorithm has been proposed [2]. Below the
application of this definition to Turing machines will
be discussed. Let (Mj ), i =0, 1, ... , be the natural
numbering of Turing machines characterized by the
fact that the machine itself (i.e. its program) may be
effectively reconstituted from its number, and the
number of the machine may be effectively reconstituted
from the machine (i.e. from its program). A general
recursive function s is a measure of the complexity of the
machines (s(i) being the complexity of the machine M j )
if and only if: a) for any y there exists not more than a
finite number of machines with complexity y; and b)
there exists an effective procedure which makes it pos-
sible to determine, for any y, all the machines with
complexity y.
Let s be an arbitrary measure of the complexity of a
Turing machine. If U is an arbitrary, effectively (i.e.
algorithmically) enumerable infinite subclass of Turing
machines, then there exists a machine T belonging to
V, and there exists a machine T' (which mayor may
not belong to V) such that T' and T compute the same
function, and the complexity of T' is much lower than
that of T. Hence it follows, in particular, that there
exist primitive recursive functions whose lowest com-
plexity in the primitive recursive form (i.e. through
schemes of primitive recursion) is much higher than
their lowest complexity in the general recursive form
(i.e. through schemes of general recursion). Let the
complexity of normal algorithms and of Turing
machines mean, respectively, the length of the encoding
and the number of internal states. Then it will be pos-
sible to realize any function of the algebra of logic in N
variables (cf. Boolean function) [3]: a) by a normal
algorithm in an m-Ietter alphabet with complexity
~2N /10& m; and b) by a Turing machine with an m-
letter external alphabet with complexity
~2N / N(m -1). The study of the complexity of algo-
rithms solving the finite restrictions of algorithmically
unsolvable problems (the so-called bounded algorithmic
problems) began in the nineteen-sixties. AA Markov
134
considered the following problem: To construct, for
any function of the algebra of logic in N variables, an
encoding of a normal algorithm in the alphabet
cP = {O, 1, a, b, c} which realizes the function and which
has least complexity among all such algorithms. It has
been shown [1], [4] that the complexity of a normal
algorithm which solves this problem is of the order 2N.
Another problem which has been studied is the com-
plexity of algorithms which solve, for the first n natural
numbers, the problem of belonging to a recursively
enumerable set (the complexity of n-segments of recur-
Sively enumerable sets). It was shown that in the case of
normal algorithms, the order of such a complexity is
not more than 10& n, and that this estimate cannot be
reduced in the general case. There also exist sets
defined by simple diagonalization of which the n-
segments have complexity order n. It was also shown
that if the time of operation of algorithms is general-
recursively bounded, then the complexity of n-segments
of recursively enumerable sets may increase exponen-
tially, and the exponent may attain value n.
Complexity of description of an algorithm is used in
formalizations of the problem of the minimum com-
plexity of an algorithm which constructs some finite
object. This minimum complexity is often named sim-
ply the complexity of the finite object (for this particular
specification of the complexity of description of an
algorithm).
A definition of the complexity of finite objects was
first proposed by AN. Kolmogorov (cf. Algorithmic
information theory). It was found that the following
asymptotically exact relations exist between the com-
plexity K(x) according to Kolmogorov, the complexity
Mm(x) of the same objects when expressed through the
length of the encoding of a normal algorithm into an
m-letter alphabet, and the complexity Tm(x) expressed
by the number of internal states of a Turing machine
with an external m-letter alphabet:
M () ...KJ& T ( ) K(x)
m x ~ lO&1 m ' m x ~ (m-l)log2 K (x)'
References
[I] MARKov, A.A.: 'Nonnal algorithms connected with the com-
putation of Boolean functions', Izv. Akad. Nauk SSSR Ser.
Mat. 31 (1967), 161-208 (in Russian).
[2] BLUM, M.: 'A machine independent theory of the complexity
of recursive functions'. J. ACM 14 (1967),322-336.
[3] KUZ'MIN, V.A.: 'Realization of functions of the algehra of
logic by automata, nonnal algorithms and Turing machines',
Problemy Kibernet. 13 (1965). 75-96 (in Russian).
[4] PETRI, N.V.: 'Algorithms connected with predicates and
Boolean functions', Soviet Math. Dokl. 10, no. 2 (1969), 294-
297. (Dokl. Akad. Nauk SSSR 185. no. I (1969). 37-39)
[5] ZVONKIN, A.K. and LEVIN, L.A.: 'The complexity of finite
objects and the developments of the concepts of information
and randomness hy means of the theory of algorithms'. Rus-
lian Math. Surveys 25, no. 6 (1970), 82-124. (Ulpekhi Mat.
Nauk 25. no. 6 (1970), 85-127)
Ya.M. Barzdin'

Editorial comments.
References
[A1] liARTMANIS. J. and HOPCROYf, J.E.: 'An overview of the
theory of computational complexity', J. ACM 18 (1971),
444-475.
AMS 1980 Subject Classification: 68C25, 68C05
ALGORITHM, COMPUTATIONAL COMPLEXITY
OF AN - A function giving a numerical estimate of the
difficulty (amount of time/storage involved) of the pro-
cess of application of an algorithm to the inputs. A
more precise definition of the computational complex-
ity of an algorithm is the concept of a cost function
(step-counting function) - defined as a decidable rela-
tion between the objects to which the algorithm is
applicable and the natural numbers, and which has a
range of definition coinciding with the range of appli-
cability of the algorithm.
Usually, the time and space characteristics of algo-
rithmic processes are considered. Thus, for a Turing
machine M the time cost function (duration of work)
TM(P) is the number t of cycles of work of M required
for the conversion of the initial configuration of Pinto
the final configuration; the memory cost function (or
space function) SM(P) is defined as the number of cells
on the tape scanned by the head of the machine. Time
and storage costs for normal algorithms (d. Nonnal
algorithm), iterative arrays, multi-head and multi-tape
Turing machines, etc., are defined in a similar manner.
The common feature of such cost functions is the
fact that there exists an effective procedure establishing
for any algorithm a (i.e. in particular for a Turing
machine or, rather, for its program), any input x and
any non-negative integer t, whether or not the process
of applying a to x will terminate with complexity t.
This observation leads to the abstract theory of compu-
tational complexity [1]. An effective procedure r is
called a computational measure if: 1) it always ends in
the result 0 or 1 when applied to triplets of the form
<algorithm, input, natural number>; and 2) it has the
property that for any algorithm a and any input x the
equality rea, x, t)= 1 is true for not more than one
natural number t, this t existing if and only if the pro-
cess of application of a to x eventually terminates. The
cost function Ra with respect to the measure r for a is
introduced if and only if rea, x, t)= 1, and then
Ra(x)=t.
This last equality is equivalent to the statement 'the
computational complexity a on x is t (in the measure
r)'.
Given some computational measure, one can con-
sider the complexity of computation of a given function
f - for example, that of finding an algorithm a which
computes f 'better than other algorithms'. However, as
ALGORITHM, COMPUTATIONAL COMPLEXITY OF AN
exemplified by the speed-up theorem (see below), such
a formulation is not always well-posed. The real prob-
lem may be the description of the rate of growth of the
cost function R a , where a computes f For example,
find upper and lower bounds for the computational
complexity of f, i.e. two functions G(x) and g(x) such
that there exists a computation a for the function f for
which Ra(x)~G(x), and such that for any algorithm /3
which computes f the function R /3 in a certain sense
majorizes g.
Another important branch of the theory of computa-
tional complexity is the study of complexity classes -
sets of functions for which there exist computations
with a complexity not exceeding some bound out of the
set of bounds which defines the class. This branch also
includes problems of comparative computational com-
plexity for various types of algorithms (automata): the
transition to another algorithmic system and another
measure of complexity is usually equivalent to consider-
ing a new measure for the first algorithmic system.
Below a few fundamental results are given which are
independent of the choice of measure of complexity
(and of the formalization), as long as there exists, for
example, an effective mutual simulation of the algo-
rithms of the type under consideration and ordinary
Turing machines. (For the sake of simplicity, it is
assumed that the subject is the computation time of
Turing machines computing functions from non-
negative integers to non-negative integers.) Let T and h
be computable integer-valued functions (d. Computable
function) on the domain of the algorithms, and let f be
a function defined on this domain which can assume
two values only - say, 0 and l.
First, there exist functions whose computation may
be as complex as one pleases. More exactly, for any
function T there exists a computable function f such
that for any algorithm a which computes f, the inequal-
ity
Ra(x) > T(x)
is false only in a finite number of cases.
Secondly, there exist functions, any computation of
which may in principle be improved as much as one
pleases. More precisely, the speed-up theorem is valid:
For any function h (e.g. h(t)=2/) it is possible to find a
computable function f such that for any algorithm a
which computes f it will always be possible to find (no
effective procedure is assumed) an algorithm /3 which
also computes f and such that for all x (except for a
finite set)
h(Rf3(x)) < Ra(x).
In case h (t) = 21 this yields
Rf3(x) < 10& Ra(x)
for almost all x.
135
ALGORITHM, COMPUTATIONAL COMPLEXITY OF AN
For the computational measures of time and space
(for Turing machines), the result of the speed-up
theorem is valid for most computations in the following
form: If f is computable with complexity R (n), it is
also computable with complexity R(n ) / 2 (one says
that f is computable with complexity R (n) if, for some
a which computes it, Ra(P)~R(n) for all words P of
length n in the alphabet under consideration). For this
reason time and space computational complexities of
concrete functions are often order of magnitude esti-
mates. Some specific results are quoted below.
The time to decide the equality of words (of length
n) is of order n 2 for Turing machines and normal algo-
rithms; any property of words which may be decided
on Turing machines in a time asymptotically less than
n lo~ n can also be decided in time n; computations by
a Turing machine of time T(n) (with the order of
T(n)~~2) can be simulated by a Turing machine with
space T(n). Furthermore, the multiplication of two
n-digit numbers on multi-tape Turing machines can be
performed in time n I +(, but not in time order n, but
iterative arrays can achieve this in real time, i.e. the i-th
digit of the product is produced before the (i + l)-th
digits of the factors are inputted. When algorithmic
processes of one type are modeled by other processes,
the computational complexity usually changes. Thus, a
reduction in the dimension of the iterative arrays
results in longer computation time. Multi-head Turing
machines and multi-tape Turing machines can simulate
each other in real time. Many types of algorithms can
be simulated on Turing machines, such that the compu-
tation time increases polynomially (usually quadrati-
cally) with an insignificant increase in space complexity
(usually, an additional constant factor).
The complexity approach is useful in studying cer-
tain hierarchies of classes of general recursive functions
(cf. General recursive function), in connection with their
logical and algebraic features. For instance, primitive
recursive functions coincide with the functions comput-
able by Turing machines in space bounded by a primi-
tive recursive function. The existence of arbitrarily
complex functions is established by diagonalization
methods. Natural analogues for such lower bounds on
computational complexity were found for the validity
problem in certain decidable theories and in areas
related to mathematical linguistics. However, for many
problems of practical importance the problem of
obtaining non-trivial lower bounds presents consider-
able difficulties. This is the case, in particular, for cer-
tain combinatorial problems which are formalized (in
one variant) as problems of finding, out of all words of
a given word length, those which satisfy a predicate
computable in polynomial time. One can also mention
the complexity of the computations performed by non-
136
deterministic automata and probabilistic automata (cf.
Automaton, probabilistic). In the non-deterministic case
the computational device can choose both ways when
presented with a choice of how to continue the compu-
tation, by 'splitting' into two independent copies, each
of which has the same capabilities as its parent. Its
computational complexity is understood to mean the
complexity of the 'simplest' successful computation
trace. In the probabilistic case, the actual course of the
computation is determined not only by the program
and the argument, but also by a random sequence of
numbers, and the result should coincide, to a great
degree of probability, with the value of the function
being computed. In both cases the computational com-
plexity can sometimes be proved to be less than that of
deterministic computations for the same problem.
The quality of algorithms is evaluated not only by
their computational complexity, but also by the com-
plexity of description an algorithm (cf. Algorithm, com-
plexity of description of an). In some publications,
both these approaches are combined.
References
[1] HARTMANIS, J. and HOPCROFT, J.E.: 'An overview of the
theory of computational complexity', J. ACM 18 (1971), 444-
475. N . V.. P etn.
Editorial comments. The area has considerably
expanded from the state described in the article above. Two
key references [A1], [A5] are listed below.
The results on Turing machines referred to, up to the
result on multiplication, hold for Turing machines in the ori-
ginal version only (single input/work tape, no extra work
tapes). In the West, the standard Turing machine is nowa-
days the (off-line) model with multiple work tapes. Turing
machines are useful, because of their simplicity, to obtain
lower bounds on the computational complexity of problems.
Upper bounds, by exhibiting actual algorithms, are usually
based on the RAM (random access machine), which more
realistically models the 'von Neumann' computer. The two
alternatives can simulate each other with a constant multipli-
cative factor in space and a polynomial time overhead. The
latter property characterizes all 'reasonable' models (First
Machine Class) for sequential computations (invariance
Thesis). The determination of such overheads is a subject of
the theory of machine models. For instance, multi-tape Tur-
ing machines with one-way input (no backup on the input
tape) can be simulated in square time by a Turing machine
with one-way input and a single work tape in an obvious
way, and it has been shown that less than square time is
insufficient [A2], [A3]. One can simulate k work tapes by 2
work tapes in time order n log n; it is not known whether this
is optimal, but the lower bound is non-linear. In contrast,
multicounter machines can be simulated in real time (step-
for-step) by a Turing machine with a single work tape [A4].
(A counter is a storage device which can hold any integer
and supports unit increments, unit decrements, and check
for zero.) The Invariance Thesis Yields a justification for the

identification of the class of problems which can be solved
deterministically in polynomial time (denoted by P) as the
class of 'tractable' problems. Adding the feature of non-
determinism gives the class NP (non-deterministic polyno-
mial time), which is known to contain a large variety of
problems which one likes to solve in daily life. The best
known deterministic algorithms for solving NP-problems
require exponential time. On the other hand, it is unknown
whether P=NP. I.e., it is unknown whether these exponen-
tial time algorithms can be replaced by deterministic polyno-
mial time ones. This question is generally considered the
most important open problem in the theory of computation.
It is known that the complexity of NP is exemplified in
specific NP problems. These NP-complete problems have
the property that P=NP if and only if such a problem can
be solved in deterministic polynomial time (is in P). Similar
questions arise with respect to the relations between dif-
ferent space complexities, and the relations between time
and space. The fundamental complexity classes related to
space are PSPACE (NPSPACE), which consist of those
problems which can be solved in polynomial space by
deterministic (non-deterministic) machines from the First
Machine Class. It is known (Savitch's theorem) that
PSPACE=NPSPACE. Clearly, PSPACE contains NP but it is
unknown whether the inclusion is strict.
There exists a large variety of models for parallel compu-
tation. Many of those models (Second Machine Class) have
the property that time on the parallel model is equivalent to
space on the sequential model, up to a polynomial overhead
(Parallel Computation Thesis). Each problem in NP can be
solved in deterministic polynomial time on a machine of the
Second Machine Class. Other models seem to be even
more powerfull, albeit totally unrealistic. The power of the
parallel models is rooted in the possible activation of at least
an exponential amount of hardware in a polynomial number
of computation cycles (called 'time'). As soon as one
enforces realistic bounds (e.g., because of the bounded
speed of light) on the amount of hardware which can be
activated by a computation in a given timespan, the Parallel
Computation Thesis becomes invalid. For many specific
problems it has been shown that they can be solved on a
parallel device with a polynomial amount of hardware in
time order log n k (NC complexity class).
References
[A 1] GARY, M.R. and JOHNSON, D.S.: Computers and intractabil-
ity; A guide to the theory of NP-completeness, Freeman,
1978.
[A2] MAASS, W.: 'Combinatorial lower bound arguments for deter-
ministic and nondeterministic Turing machines', Trans. Amer.
Math. Soc. 292 (1985), 675-693.
[A3] VITANYl, P.M.B. and LI, M.: 'Tape versus stack and queue:
the lower bounds', Information and Computation (To
appear).
[A4] VITANYl, P.M.B.: 'An optimal simulation of counter
machines', SIAM J. on Computing 14 (1985),1-33.
[A5] WAGNER, K. and WECHSUNG, G.: Complexity theory, Reidel,
1986.
AMS 1980 Subject Classification: 68C05, 68C25
ALGORITHM IN AN ALPHABET, algorithm over an
ALGORITHM, LOCAL
alphabet A - 'An exact, generally understandable set of
instructions describing a potentially realizable process
of successive transformation of abstract words over the
alphabet A, in which any word over A may serve as the
initial word' [1].
An algorithm over an alphabet is a special case of
the general concept of an algorithm. The input and the
possible outputs of an algorithm over an alphabet are
constructive objects (cf. Constructive object) of suffi-
ciently general type - words (cf. Word). Various for-
malizations of an algorithm over an alphabet were
developed in the nineteen-seventies; the best known
ones are Markov's normal algorithms (d. Normal algo-
rithm) and algorithms defined using the Turing
machine. When constructing specific algorithms over
alphabets, it is often useful to extend the starting
alphabets by introducing additional letters and consid-
ering the algorithms over the alphabets thus extended.
In the case of two algorithms over the same alphabet,
the concepts of equivalence and complete equivalence
of these algorithms with respect to it can be naturally
introduced.
References
[1] MARKov, A.A.: Theory of algorithms, Israel Progr. Sci. Trans!.,
Jerusalem, 1961 (translated from the Russian). Also: Trudy
Mat. Inst. Steklov. 42 (1954).
N.M. Nagornyf
Editorial comments. The connection between algo-
rithmic descriptions of languages (on the basis of algorithms
over an alphabet) and generative descriptions using formal
grammars in the style of N. Chomsky, constitutes the
essence of the theory of abstract automata and formal
languages. The classic text is J.E. Hopcroft and J.D.
Ullman's book [A1].
References
[A1] HOPCROFT, J.E. and ULLMAN, J.D.: Formal languages and
their relation to automata, Addison-Wesley, 1969. Revised
version: 'Introduction to automata, languages and computa-
tion' Addison-Wesley, 1979.
AMS 1980 Subject Classification: 68C05, 03003,
03010
ALGORITHM, LOCAL - An algorithm which speci-
fies the attributes of elements of a structured set and
uses, at each step, only information about neighbour-
hoods of the element. Problems of the existence or
non-existence of effective algorithms for various
discrete optimization problems are naturally formulated
in terms of local algorithms.
Let there be given a family {WC: WCEM} of sets, and
let there be assigned to each pair (~, WC), where ~ E WC,
a set S(~, WC) such that: 1) ~ES(~, WC); 2)
S(~, WC)~WC; and 3) if ~EWCl' ~EWC2 and
S(~, WCl)~WC2 ~WCl' then S(~, WCl)=S(~, WC 2). The
set S(~, WC) will then be denoted as a neighbourhood of
~ in WC. Examples of neighbourhoods are conjunction
137
ALGORITHM, LOCAL
neighbourhoods in a contracted normal form (cf.
Boolean functions, minimization of).
Let f be a finite undirected graph, f =M 1 U M 2,
where M 1 = {a I> ,aq } is the set of vertices of f,
and let M 2 ={(aj" aj), ... ,(a", a,,)} be the set of
edges of f. A neighbourhood SI(R, f) of an edge
R =(aj, aj) of the graph f consists of all vertices of the
edges incident with the edge R, and also of all edges
whose vertices are included in the neighbourhood
S 1(R, f). Let Sn -1 (R, f) be a defined neighbourhood;
the neighbourhood Sn(R, f) then consists of all vertices
of the edges incident with the vertices of the neighbour-
hood Sn -1 (R, f) and all edges, the vertices of which
belong to the neighbourhood Sn(R, f). The neighbour-
hoods SI(aj, f), ... ,Sn(aj' f) of an arbitrary vertex aj
of the graph f are introduced in a similar manner.
Let a system of two-place predicates
P 1 (~, 9R), ... ,PI(~' 9R), which is subdivided into two
non-intersecting subsets <P" ... ,P, >
<Pr+1> ... ,PI>, be defined on pairs (~, 9R), ~E9R.
The elements of the first subset are called the main
predicates, while those of the second are called the auxi-
liary predicates. A vector a=(ab ... ,al) is called an
information vector if aj E {O, 1, 6. }, i = 1, ... ,I. A vector
a is called permissible for ~ in 9R if for all aj=l=A the
equality aj = Pj(~, 9R) is satisfied. The set J (~, 9R) of
all information vectors permissible for ~ in 9R is called
the information set of ~ in 9R.
Let 9R= {~b ... ,~q} and (ail, ... ,ai/)EJ (~j, 9R),
.I -- 1, , q. The set on* - {ora" ... all
;U~ - '?t}
oraq' ... a" }
, ,u q
is called permissible for 9R. The class J(9R) of all sets
9R* permissible for 9R is called the information class of
the set 9R over the system of predicates P" ... ,PI'
Clearly, a neighbourhood S(~, 9R) defines a neighbour-
hood S(~a, ... a" Wt).
Introduce a system of functions cp], ... ,CPI such that
cM~a, ... a" S(~a, ... a" [If) = (f3l. ... ,f3I)'
The functions cPj will be defined over all pairs
(~a, ... a" S(~a, ... a" 9R* such that ~I' a, E9R*,
9R* EJ(9R), and satisfies the fo~~wing conditions: 1)
aj =/3j if j=l=i; and 2) the set 9R , which is obtained
from 9R* by replacing the e~etpent ~a, ... a, by ~.81 .. '.8',
is permissible for 9R, i.e. 9R EJ(9R). For the sake of
brevity, the pairs (~I' .. a" S(~a, ... a" 9R*) are
denoted by(~,a], ... ,ai'S, 9R*).
A partial order is introduced on some of the above
sets: 1) on the set M] ={O, 1, A} one takes A<O,
A< 1; 2) on the set M 2 of information vectors of length
lone takes (a], ... ,aa~(/3], ... ,/31) if a.i~/3i'
i = 1, ... ,I; 3) on the set of elements wi~h marks one
takes ~a ... a'~~.8 ... .8'1'f (a], ... , al )~(/3
I I """ ], ... , /3I;
4) on the set M* = U J(9R) one takes 9R~ ~9R; if,
* * IJJIcM
first, [)C] and [)C2 belong to the same information class
138
J(9R) and, secondly, if (al> ... ,al)~(fJ}, ... ,PI) fol-
lows from ~, '" a, E9Ri and ~P, ... P, E9Ri; 5) on the
set of neighbourhoods of the form S(~' ... a" 9R*) one
takes SI ~S2' where SI =S(~, .. 'a" 9Ri) and
S2 =S(~p, ... P" 9Ri) if and only if
S(~, 9R])=S(~, 9R2) and if it follows from
?BY' ... y, ES 1 and ?B6, ... 6, ES 2 that
(y}, ... 'Y/)~(8}, ... ,81),
Let A and B be elements of one of the sets 1) - 5). If
A ~B and A ~B, then the elements A ap.d B are called
equi-informative, which is denoted by A B. A function A
cp(~, aI, ... ,ai, S, 9R*) is called monotone if it follows
from SI ~S2 that, for i= 1, ... ,I,
cM~, IXl. .. ,IXl. S l. 9nj) .;;;; cMI}!, f3l. . ,f3l. S 2, '1n 2).
The determination of a local algorithm also necessi-
tates the introduction of an ordering algorithm A '/I' Let
M be an arbitrary set of elements with informatio!l vec-
tors, and let N ={I, ... ,I}. Consider the set M X N of
and
all pairs (~I'" 0/, ,j) such that ~" .. 0/, E, j EN,
aj=6.. The algorithm A'/I orders the set MXN. The
local algortihm A is completely defined by the system
of predicates PI>' .. ,PI, by the subdivision of this sys-
tem into main predicates P" ... , Pr and auxiliary
predicates Pr +" ... ,PI, by the system of monotone
functions cP" ... ,CP/, cPj =CPj(\}J, a}, ... ,ai> S, 9R*), and
by the algorithm A '/I'
Let [)C* = U~]~' ... ai" 9R* EJ(9R). The first step
of the local algorithm may be described ..as follows.
Algorithm A'/I is applied to the set M X N, where
M = 9R*; for the first pair (~' ... a, ,j) of the ordered
sequence one computes
cpi~, a], ... ,ai, S, 9R*)=(fJI' ... ,PI) and the element
~, 0/, is then replaced by \}JP' ... .8, ; if
(ab ... ,al)=(fJb ... ,PI)' then one takes the second
pair, etc.; if for all elements (?BY' ... Y', i) the equality
cp(?B, yJ, ... ,YI, S, [)C*)=(y], ... ,YI) is true, then the
local algorithm A is terminated after reviewing all the
pairs from M X N. Otherwise, after replacement of the
vector (ab'" ,al) by the new vector (/3], ... ,/31)'
algorithm A is terminated if there are no more elements
with at least one symbol A in the first r places of the
information vectors; if there are such elements, the first
step of the local algorithm terminates. Let the number
of steps of the local algorithm A which have been per-
formed be n. The (n + 1)-st step is carried out exactly as
the first step, except that instead of the set [)C* one now
considers the set [)C~ into which the set [)C* has been
converted after the first n steps of the local algorithm
A. Owing to the monotone nature of the functions CPi,
i = 1, ... ,I, the local algorithm A will be terminated
) after a finite number of steps.
The starting theorems for local algorithms are the
uniqueness theorem and the theorem on the existence of a

best algorithm. The first theorem states that the result
of the calculation of the main predicates is independent
of the algorithm A" (the order in which the elements of
the set 9JC' are taken). The second theorem stipulates
the existence, under very general conditions, of a best
local algorithm, i.e. of a local algorithm which uses a
given system of neighbourhoods to compute given main
predicates with fixed auxiliary predicates, whenever this
is done by any other local algorithm with an identical
system of neighbourhoods, and main and auxiliary
predicates.
The problem of finding a best local algorithm in
explicit form has been solved for algorithms of syn-
thesis of minimal coverings. Let there be given a tuple
<~I"" ,~q,/1.I"" ,}1.q,M>, where m: i is a set, Jl.i is
the complexity of m:i, }1.i>O, i=l, ... ,q,
M c:; U r= I m:i . The complexity of the tuple
(m:il , . . . , ~i) is ~~ = I Jl.i The problem is how to con-
j References
struct a covering of M by sets from a number of
m: 1, ,~q with minimum complexity. A system of
neighbourhoods SI, ... ,Sn, ... , for each ~i is intro-
duced in a manner similar to the introduction of a sys-
tem of neighbourhoods for conjunctions. A best local
algorithm is constructed for a system of neighbour-
hoods Sk at fixed predicates: the set of auxiliary predi-
cates is empty; the main predicates are considered to
be the predicates PI (m:, M) (m: is not included in any
minimal covering of 9JC) and P 2 (m:, M) (m: is included
in all the minimal coverings of 9JC). Best local algo-
rithms have also been constructed to calculate simple
properties of edges of a graph. Various local algorithms
have been proposed to solve problems of minimization
of Boolean functions, of discrete linear programming,
etc.
Computability of best local algorithms in a class of
local algorithms is an important problem in the theory
of local algorithms. If a system of imbedded neighbour-
hoods S I (m:, 9JC) c:; ... c:; Sn(m:, 9JC) c:; . .. is given on
the pairs (~, 9JC), ~E9JC, 9JCEM, and if the local algo-
rithm operates with respect to the system S(~, 9JC) of
neighbourhoods where
Sk-I(m:, 9JC)C:;S(m:, 9JC)C:;Sk(m:, 9JC), then one says that
the local algorithm has index k. It is natural to consider
the number of predicates PI, ... ,PI in the definition
of a local algorithm as the quantity of memory of the
local algorithm. Let there be given a set (JJ of predi-
cates P(m:, 9JC). It is considered that all the predicates
participating in the definition of the local algorithm
belong to (JJ. Let A(9JC') be the result of applying the
local algorithm A to 9JC', 9JC' EJ(9JC), 9JCEM. One says
that the predicate P I (~, 9JC) is not (k, I)-computable if
for any local algorithm of index k with memory quan-
tity I for the main predicate P I and for the auxiliary
predicates P 2 , . ,PI from (JJ there exists a set 9JC'
ALGORITHM, REPRESENTATION OF AN
such that in A(9JC') the predicate PI will not be com-
putable for all elements.
Let f(x I, . . . ,xn ) be the set of all Boolean functions
(cf. Boolean function) in the variables x I, . . . ,Xn and
let P I (~,f (x I, . . . ,xn be the predicate ('the conjunc-
tion ~ is included in at least one minimal disjunctive
normal form'). With natural limitations imposed on the
class of predicates (JJ it was found that the predicate
P I (~,f (x I, . . . ,xn is not (k, I)-computable if
kl<const2n It is interesting to note that the predicate
P(m:,f(XI, ... ,xn ('the conjunction ~ is included in
at least one dead-end disjunctive normal form of 1')
(cf. Boolean functions, nonnal fonns of) is (2, 1)-
computable, but is not (1, I)-computable. Similar
results were obtained for predicates describing the entry
of an edge into a shortest path between given vertices
of a graph.
[IA] ZHURAVLEV, YU.I.: 'Local algorithms for the calculation of
infonnation Part 1', Cybernetics 1, no. I (1965), 10-16. (Kiber.
netika (Kiev) 1, no. I (1965), 12-19)
[lB] ZHURAVLEV, Yu.I.: 'Local algorithms for the calculation of
infonnation Part II', Cybernetics 2, no. 2 (1966), 1-8. (Kiber
netika (Kiev) 2, no. 2 (1966), I-II)
[2] KHUTORYANSKAYA, LV.: 'Aspects of the theory of local algo-
rithms on graphs', Cybernetics 7, no. I (1971),37-44. (Kiber
netika (Kiev), no. I (1971), 29-33)
[3] ZHVRAVLEV, Yu.I.: 'Set-theoretical methods in the algebra of
logic', Problemy Kibernet. 8 (1962), 5-44 (in Russian).
See also the references to Boolean functions, minimiza-
tion of.
Yu.I. Zhuravlev
Editorial comments. As a specific area (or in this form),
the subject seems not to have been studied in the Western
literature. It may be related to the study of Systolic Array
algorithms [A1], Chapt. 8.
References
[A 1] MEAD, C. and CONWAY, L.: Introduction to VLSI systems,
Addison-Wesley, 1980, pp. 263-330.
AMS 1980 Subject Classification: 68C05
ALGORITHM, REPRESENTATION OF AN, coding
of an algorithm - A constructive object of a certain kind
(as a rule, a natural number or a word) containing
complete information about the algorithm, coded in
accordance with the rules laid down for this type of
algorithm. A coding of an algorithm is usually so
defined that the procedures for obtaining the coding of
the original algorithm and obtaining the original algo-
rithm from its coding are as simple as possible.
The coding ~I of a nonnal algorithm ~ over an
alphabet A (not containing the letters a, {3 and y) is
defined as a word over the alphabet A a{3y which is
obtained as follows [1]. In each substitution formula of
the scheme of ~ the arrow is replaced by the letter a,
the point (if any) by the letter {3, and the letter y is
written at the end of the word thus formed. The result-
139
ALGORITHM, REPRESENT AnON OF AN
ing words are then written out, one after the other, in
the same sequence as that of their respective substitu-
tion formulas in the scheme of 2{. In actual fact, 2{I is
merely a scheme of the normal algorithm 2{, written in
a somewhat different form - viz. as a word. For rea-
sons connected with the technical details of the defini-
tion of a normal algorithm, another type of coding of a
normal algorithm is somewhat more convenient - the
so-called record of a normal algorithm [1], which is the
result of translating the words of 2{I into some two-
letter alphabet (usually {O, I}). The coding of the pro-
gram of a Turing machine can be carried out in a simi-
lar manner. The role of the coding of a recursive func-
tion is played by the GOdel number of the system of
equations which define this function [2].
In any formalization of the general concept of an
algorithm, its coding is so defined that it (in fact, the
algorithm coded by it) might be included in the inputs
for the algorithms of this type. This makes it possible
to demonstrate the so-called 'theorems on universal
algorithms', dealing with the realizability of algorithms
whose codings model the operation of any algorithm of
this type.
If a coding of an algorithm is a permissible input for
the algorithm, this algorithm is called self-applicable if it
is applicable to its own coding, and not self-applicable
otherwise. It is possible to show, with the aid of a vari-
ant of the Russell paradox (cf. Antinomy), which is
applicable to this situation, that the algorithmic prob-
lem which then arises - to wit, how to recognize self-
applicable algorithms among other algorithms of this
type - is not solvable by an algorithm of this type.
This result forms the base of numerous theorems on the
un solvability of algorithmic problems.
The length of the code of an (encoding of an) algo-
rithm is a natural measure of the complexity of an
algorithm (the amount of memory which is required to
remember the program of the algorithm) (cf. Algorithm,
computational complexity of an; Algorithm, complexity
of description of an).
References
[1] MARKov, A.A.: Theory of algorithms, Israel Progr. Sci. Trans!.,
1961 (translated from the Russian). Also: Trudy Mat. Inst.
Steklov. 42 (1954).
[2] KLEENE, S.c.: Introduction to metamathematics, North-Holland
& Noordhoff, 1951.
N.M. Nagornyr
AMS 1980 Subject Classification: 03D05, 68C05
ALGORITHMIC INFORMATION THEORY - The
branch of mathematical logic which gives a more exact
formulation of the basic concepts of the theory of
information, based on the concepts of an algorithm and
a computable function. Algorithmic information theory
attempts to give a base to these concepts without
140
recourse to probability theory, so that the concepts of
entropy and quantity of information might be applica-
ble to individual objects. It also gives rise to its own
problems, which are related to the study of the entropy
of specific individual objects.
The key concept of algorithmic information theory is
that of the entropy of an individual object, also called the
(Kolmogorov) complexity of the object. The intuitive
meaning of this concept is the minimum amount of
information needed to reconstruct a given object. An
exact definition of the concept of complexity of an
individual object, and (on the basis of this concept) of
the concept of the quantity of information in an indivi-
dual object y relative to an individual object x, was
given by A.N. Kolmogorov in 1962 - 1965, after which
the development of algorithmic information theory
began.
Without loss of generality, only binary words (linear
strings of O's and l's) are considered. The complexity of
a word x is understood to mean the length I(p) of the
shortest binary word p which contains all the informa-
tion required to reconstruct x with the aid of some
definite decoding procedure. A precise definition of this
concept follows.
Let F(s) be an arbitrary universal partial recursive
function. Then the complexity of the word x with
respect to F will be
min/(p) if F(P)=x,
KFeX) = { 00 if there is no p such that F(P)=x.
The word p such that F(P)=x is known as the code,
or the program, by which F reconstructs the word x.
The following theorem is valid: There exists a partial
recursive function F 0 (an optimal function), such that
for any partial recursive function F the following ine-
quality holds:
where CF is a constant which is independent of x. This
theorem yields an invariant (up to an additive constant)
definition of complexity: The complexity K(x) of the
word x is the complexity KFo (x) with respect to some
partial recursive function F 0, which has been fixed
once and for all. Clearly, K(x)~/(x)+C, where C is a
constant independent of x.
Using K(x) one may also determine the complexity
K (x, y) of pairs of words (x, y); to do this, the pairs
(x, y) are effectively coded as one word c(x, y), and the
complexity K(x,y) is understood to mean K(c(x,y.
P. Martin-LOf studied the behaviour of the complex-
ity of the initial segments of infinite binary sequences.
Let Wn denote the initial segment of a sequence w, con-
sisting of the n initial characters. Then almost all
(according to the Lebesgue measure) infinite binary
sequences W have the following properties: a) for infin-

itely many natural numbers n the inequality
K(wn)';;;:'n -c, where c is some constant, is valid; and b)
for all natural numbers n larger than some value n.,
K(wn);;;:'n -(1 +)lo~ n, where is an arbitrary positive
number (the Martin-LOf theorem). Conversely, for any
sequence W there exist infinitely many n such that
K(wn)~n -lo~ n.
The concept of complexity is also employed in the
study of algorithmic problems. Here, a more natural
concept is the so-called complexity of solution of the
word x. The intuitive meaning of this concept is the
minimum quantity of information required to find, for
each number i~l(x), the i-th character of word x (the
length of the word x need not necessarily be known).
More exactly, the complexity of solution of the word
x =x} ... Xl(x) by a partial recursive function G(s, t)
means that
min{I(P): Vi""-/(x), G(p, i)=x;},
KRG(x) = {
ereIS no sueh p.
00 if t h
The following theorem is valid: There exists a partial
recursive function Go (an optimal function) such that
any other partial recursive function G satisfies the ine-
quality KRG.(x )~KRG(x) + CG, where CG is a constant
independent of x. The complexity of solution, KR(x),
of the word x is the complexity KR G, (x) of some
optimal partial recursive function Go which has been
fixed once and for all. Clearly, KR(x)~K(x)+C and
K(x)~KR(x)+21(x)+C. It is possible, by using
KR(x), to determine, for any set M of natural numbers,
the complexity K(M, n) for the n-segment of the set M:
K(M, n)=KR(wn), where w=x} Xi . .. is the
characteristic sequence of the set M (x; = 1 if i EM, and
Xi =0 if HiM).
Algorithmic problems are usually represented as
problems of belonging to some recursively enumerable
set M. If some value n is given, and the problem of
belonging to M is restricted to the first n natural
numbers only, then one obtains a bounded algorithmic
problem. The magnitude K(M, n) may be intuitively
understood as expressing the quantity of information
required to solve this bounded problem. In particular,
the set M is recursive if and only if K (M, n) is
bounded from above by some constant.
It follows from the Martin-Lof theorem that there
exist sets for which K(M, n)-n. It was also found that
sets with maximum degree of complexity exist even
among sets defined by arithmetic predicates with two
quantifiers. However, the following theorem applies to
recursively enumerable sets: 1) for any recursively enu-
merable set M and any value n, the inequality
K(M, n )~lo~ n + C, where C does not depend on n, is
valid; and 2) there exists a recursively enumerable set
M such that for any n the inequality K(M,nlo~ n
ALGORITHMIC INFORMATION THEORY
holds. At the same time there exists a recursively enu-
merable set M such that if the time t of computation of
an arbitrary general recursive function is bounded, the
estimate Kt(M, n);;;:.n / c" where Ct does not depend on
n, is valid.
Sets which are universal with respect to a certain
type of reducibility may also be defined in these terms
(cf. Universal set): A set M is weakly tabularly univer-
sal if and only if there exists an unbounded general
recursive function f (n) such that the inequality
K(M, n);;;:. f (n) is valid for all values of n.
In studying the complexity of bounded algorithmic
problems other measures of complexity may also be
employed, e.g. the minimum of the length of the coding
of the nonnal algorithm which solves the bounded prob-
lem in question. It was found, however, that asymptoti-
cally exact relations (d. Algorithm, complexity of
description of an) exist between the complexities intro-
duced above and the analogues of these complexities
expressed as the minimum length of the coding of an
algorithm (or as the minimum number of internal states
of a Turing machine).
Another concept which is introduced in constructing
algorithmic information theory is that of conditional
compleXity of a word x, for a known y, using a partial
recursive function G(s, t):
min {l(P): G(p,y)=x},
KG(x Iy) = { 00 if there is no p such that G(p,y)=x.
This concept is related to the theorem on the existence
of an optimal function Go, and the conditional com-
plexity K(x Iy) of a word x, if y is known, is defined as
the complexity KG, (x Iy) using an optimal function Go
fixed once and for all. The intuitive meaning of
K(x Iy) is the minimum quantity of information which
must be added to the information contained in the
word y for the reconstitution of the word x to be possi-
ble. Clearly, K(x ly)~K(x)+C.
The next central concept of algorithmic information
theory is the concept of the amount of information in an
individual object y relative to an individual object x:
I(y:x) = K(x)-K(x Iy).
The quantity /(y :x) is called the algorithmic amount of
information in y about x. The quantities K(x) and
K(x Iy) are sometimes called the algorithmic entropy of
x and the algorithmic conditional entropy of x for a
given y. The formulas for the decomposition of the
entropy H(X, y)=H(X)+H(YIX) and of the commu-
tation of information /(X:Y)=/(Y:X) are valid in the
algorithmic concept up to terms of order
O(1ogH(X, Y:
I K(x,y)-[K(x)+ K(y Ix)ll ""- o (log K(x, y),
I I(x:y)-I(y:x) I ""- Q(logK(x,y)).
141
ALGORITHMIC INFORMATION THEORY
There exists a certain relation (due to Kolmogorov)
between the algorithmic and the classical definitions of
the amount of information (or, more exactly, between
the complexity of a word according to Kolmogorov and
the entropy of the Shannon distribution of frequen-
cies): Let a number r be given, let the word x of length
i'r consist of i words of length r, and let the k-th word
of length r appear in x with a frequency qk
(k = 1, ... ,2'); one then has
K(.x) ..;; H +a:(i),
I ALGORITHMIC LANGUAGE, formal programming
where
t
H = - L qk 10& qb
k=l
and a(i)~O as i~oo. In the general case no closer con-
nection can be established between entropy and com-
plexity. This is because entropy is intended to study
sequences subject to frequency laws only. Conse-
quently, one may establish the following relation for
random sequences w in measure between the magni-
tudes under study, which corresponds to independent
trials:
lim K(~ir) = H.
1---"00 I
Similar considerations apply to an arbitrary ergodic
stationary stochastic process.
References
[1] KOLMOGOROV, A.N.: 'Three approaches to the quantitative
definition of information', Problems Inform. Transmis-
sion 1 (1965), 1-7. (Problemy Peredachi Informatsii 1, no. I
(1965),3-11)
[2] KOLMOGOROV, A.N.: 'On the logical foundations of informa-
tion theory and probability theory', Problems Inform.
Transmission 5, no. 3 (1969), 1-4. (Problemy Peredachi Infor-
matsij 5, no. 3 (1969), 3-7)
[3] ZVONKIN, A.K. and LEVIN, L.A.: 'The complexity of finite
objects and the developments of the concepts of information
and randomness by means of the theory of algorithms', Rus-
sian Math. Surveys 25, no. 6 (1970), 82-124. (Uspekhi Mat.
Nauk 25, no. 6 (1970),85-127)
[4] BARDZIN', YA.M.: 'Complexity of programs to determine
whether natural numbers not greater than n belong to a recur-
sively enumerable set', Soviet Math. Dokl. 9, no. 5 (1968),
1251-1254. (Dokl. Akad. Nauk SSSR 182, no. 6 (1968),1249-
1252)
[5] KANOVICH, M.l.: 'Complexity of resolution of a recursively
enumerable set as a criterion of its universality', Soviet Math.
Dokl. 11, no. 5 (1970), 1224-1228. (Dokl. Akad. Nauk SSSR
194, no. 3 (1970), 500-503) v M B d"
1a. . arz
Editorial comments. The theory of Kolmogorov com-
plexity, or 'algorithmic information theory', originated in the
independent work of R.J. Solomonoff, Kolmogorov, L. Levin,
Martin-Lof and G.J. Chaitin. Recently, J. Hartmanis has
refined the theory by also considering the amount of work
(time complexity) involved in reconstructing the original data
from its description.
References
[A 1] HARTMANIS, J.: 'Generalized Kolmogorov complexity and the
structure of feasible computations', in Proc. 24th IEEE Sym-
142
posium on Foundations of Computer SCience, 1983, pp.
439-445.
[A2] SoLOMONOFF, R.J.: 'A formal theory of inductive inference,
Part I and Part II', Information and Control 7 (1964), 1-22.
[A3] MARTIN-LOF, P.: 'The definition of random sequences', Infor-
mation and Control 9 (1966), 602-619.
[A4] CIiArrIN, G.J.: 'Randomness and mathematical proof', Scien-
tific American May (1975),47-52.
[A5] CHAiTIN, G.J.: 'Algorithmic information theory', IBM J. Res.
Oev. 21 (1977), 350-359.
AMS 1980 Subject Classification: 68C25, 94A 17
language - A fonnal language intended for describing
computational processes or, equivalently, for writing
down algorithms to be executed by computers. One dis-
tinguishes between problem-oriented algorithmic
languages (high-level languages), which are not related
to any specific machine, and machine-oriented algo-
rithmic languages (low-level languages), which take the
specific features of a given machine into account
(instruction set, addressing modes, etc.). The term
'algorithmic language' usually refers to a problem-
oriented language, as opposed to machine code, which is
a notation that is directly interpreted by a machine.
For the well-formed texts of an algorithmic language
(programs, cf. Program) a general algorithm defines
their execution in a unique way, which is the distinc-
tion between algorithmic languages and non-algorithmic
programming languages, for which the execution process
for a text is fully undetermined or the text merely
serves as material for the synthesis of an algorithm.
As in natural languages, an algorithmic language is
constructed over an alphabet of basic symbols (in
which the program is written down) in the form of a
hierarchical system of grammatical elements, between
which relations are given (similarly to the words,
phrases and sentences in a natural language, whose
connections are given by syntactic rules). The lowest
level elements, formed by chains of basic symbols, are
called lexemes, or lexical units. For the lexemes occur-
ring in a program the class to which they belong is
defined, and for certain classes of lexemes (e.g. identif-
iers) also their scope - some uniquely identifiable part
of the program to which all occurrences of a given lex-
erne belong (a block). Exactly one occurrence of such a
In
lexerne is said to be defining; the other occurrences of
the lexeme in its scope are called applied.
The further levels of elements of an algorithmic
language are formed by notions, or non-terminals. The
relation that may hold between the notions of an algo-
rithmic language is that of being a (direct) constituent
of (i.e. an immediate constituing part), while the indivi-
dual constituents of a given notion are related to each
other by concatenation (textual sequence). The transitive
closure of the constituent relation uniquely assigns to

each notion some subword of the text of the program,
which is said to be the (terminal) production of this
notion. There is one initial notion, the production of
which is the entire program text. A tree whose root is
the initial notion, whose terminal vertices (leafs) are
lexemes and basic symbols, whose internal vertices are
concepts and whose branches are constituent relations,
is called a production or syntax tree of a program. The
construction of such a tree is known as the syntactic
analysis or parsing of a program.
Notions and lexemes have attributes, i.e. certain sets
fixed by the description of the algorithmic language.
Determining the attributes of the elements occurring in
a program is called its semantic analysis. Finding the
attributes of lexemes begins with the analysis of their
defining occurrences, which usually contain explicit
information about the attributes. The attribute informa-
tion is transferred to all the applied occurrences of the
lexeme within the corresponding scope (identification).
The attributes of a certain notion are found, by induc-
tion over the production tree, as a function of the attri-
butes of its constituents. A substantial part of the
semantic analysis, which is very valuable for checking
the correctness of a program, is a check on the compa-
tibility of the attributes.
The rules of syntactic analysis are laid down by the
generative grammar of the algorithmic language (d.
Grammar, generative) or by an analyzing automaton
(more precisely, by its different generalizations) for
converting the program text into a production tree. The
rules of semantic analysis are usually described infor-
mally, but attempts have been made to formalize the
definition of attribute information and to take the con-
text into consideration with the aid of the mechanism
of two-level grammars (cf. AJgol-68).
The execution algorithm of the programs of an algo-
rithmic language is usually given in the form of a
correspondence between the vertices of the production
tree and various execution procedures (also called
transducers or semantic procedures). Each procedure
performs certain operations as a function of the attri-
butes and of some context of the given notion (or lex-
erne), and then determines the next procedure. The
operations of the procedure may specify the computa-
tion process directly in terms of some abstract machine
(semantics of the interpretation type), or may transform
the production of a given notion into a fragment of the
program in some high-level language (translational
type).
In specific algorithmic languages, the alphabet of
basic symbols usually consists of the letters of the
Roman alphabet (which may be replaced or augmented
by characters from a national set), digits, pairs of de-
limiters (parentheses), separators (punctuation marks)
ALGORITHMIC LANGUAGE
and symbols for certain operations. The principal
classes of lexemes are numerals to represent numbers,
literals to represent texts, and identifiers to denote the
various objects of the program that are defined in the
program itself. The principal such objects are variables,
labels (which name different parts of the program) and
procedures (which denote functions). The meanings and
the names of certain identifiers are defined by the
description of an algorithmic language (reserved
words).
Among the concepts of an algorithmic language are
basic structures - definitions, expressions and opera-
tors. Descriptions are sources of attribute information
which is assigned to the defining occurrence of the lex-
erne. Essentially, attributes characterize the type of
values calculated by executing the program, their
representation and the way in which they are stored in
the memory of the computer. For composite values like
arrays (vectors, matrices) and structured values
(records) the mode of access to their elementary com-
ponents is also specified. Expressions are the source of
values; operators are units of finite operations in the
program; basic operators are the operators of assignment
of the value of an expression to a variable, the operator
of control transfer (unconditional or conditional jump,
or goto), the operator of procedure call, and the opera-
tor of a loop.
Loops and procedures are the most characteristic
tools of an algorithmic language, providing as they do a
concise notation for very long computations. The loop
operator prescribes a repetitive execution of some part
of the program (the body of the loop), the number of
repetitions being determined by some given condition.
The definition of a procedure introduces an abbrevi-
ated, parametrized notation for the function of some
part of the program (the procedure body). The execu-
tion of the procedure body may be subsequently
treated as an elementary operation triggered by a call
of the procedure, at which time the actual values are
assigned to the parameters included in the procedure.
An important property of an algorithmic language is
its capacity to structure a program text, which makes it
easier to write, to learn and to check. The principal
means of structuring are limitations on the structure of
scopes (block structure), on procedure and loop, bodies
and rules imposed on the use of goto operators, which
limits the number of branches on the execution paths.
Programming a computer in an algorithmic language
necessitates the availability of special software in the
form of programming processors (compilers, interpreters)
- which are intermediaries between the program in the
algorithmic language and the machine. A complete pro-
cessor executes the following functions: input of the
program, lexical analysis (isolation and classification of
143
ALGORITHMIC LANGUAGE
lexemes), syntactic and semantic analysis, with indica-
tions of formal errors, synthesis of an intermediate
form (representation of the program in the internal
language of some abstract computer, which is con-
venient for the subsequent processing or execution of
the program), optimization (systematic transformations
of the intermediate forms which improves certain
characteristics of the program such as its size, speed
and demands on memory), code generation (construc-
tion of the machine program) and execution of the pro-
gram. In processors of the translational type (transla-
tors, compilers, program generators) the execution of the
program takes place after the completion of the con-
struction of the machine program. In processors of the
interpretative type (incremental translators, interactive
translators) the program is executed with the aid of
some mechanism of interpretation of either its inter-
mediate form, or of the production tree or even of the
initial text.
One distinguishes between one-phase, two-phase and
multi-phase translation schemes of programs. In the
one-phase scheme all functions of the translator are
executed during a single pass of the text of the pro-
gram. There is no optimization and no intermediate
form, while the isolated vertices of the production tree
immediately generate machine commands of the pro-
gram. In two-phase translators, the production tree is
usually built during the first pass of the program, while
the semantic analysis, combined with the generation, is
effected during the second pass. Multi-phase schemes,
involving the utilization of an intermediate form, are
employed in optimizing translators or in systems with
mixed generation for various types of machines; in cer-
tain very complex algorithmic languages syntactic and
semantic analysis is also required.
The number of algorithmic languages which may be
employed with computers is very large (more than one
thousand), but only a few of them are extensively used.
These include Algol; Algol-68; Cobol; Lisp; PL/I;
Simula; Fortran; and in the USSR also Algams; Al'fa
and Refal.
References
[1] INGERMAN,P.: A syntax-oriented translator, Academic Press,
1966.
[2] Programming languages, Moscow, 1972 (in Russian; translated
from the English).
[3] HOPGOOD, F.R.A.: Compilation techniques, American Elsevier,
New York, 1969.
[4] HiGMAN, B.: A comparitive study of programming languages,
American Elsevier, New York, 1968.
A.P. Ershov
Editorial comments. In English-speaking computer sci-
ence community, usually no specific distinction is made
between the concepts of 'progamming language' and
'algorithmic language', although the latter term historically
refers mainly to languages of the Algol family. The termi-
144
nology and concepts of. the theory of programming
languages are not standardized; the article above leans
heavily on the somewhat idiomatic approach employed in
the definition of Algol-68.
A language that has gained widespread use since the
nineteen-seventies is Pascal, which is a close relative of
Algol-60 offering some additional facilities. Next to the class
of procedural (or operational) languages, which basically
prescribe the sequence of operations to be performed, there
are functional (or applicative) languages, in which a pro-
gram is basically a collection of recursive function defini-
tions, and deductive (or logic-programming) languages, for
which a program consists of a collection of recursive defini-
tions of predicates. Most existing algorithmic langages are
procedural. An example of a functional language is Lisp,
whereas Prolog is a deductive language.
References
[A1] PRArr, T.W.: Programming languages: design and implemen-
tation, Prentice Hall, 1975.
AMS 1980 Subject Classification: 68BXX
ALGORITIlMIC PROBLEM - The problem of find-
ing a (unique) method (an algorithm) to solve an infin-
ite series of individual problems of the same type.
Algorithmic problems arose and were solved in various
branches of mathematics throughout its history; how-
ever, some of them could not be solved for a long time.
The reason for this only became apparent in the
nineteen-thirties, when an exact definition of an algo-
rithm was given in mathematical logic. It was found
that algorithmic problems can be unsolvable, i.e. that
the algorithm sought need not exist at all. As a result,
the approach to algorithms on the part of mathemati-
cians underwent a radical change, and new algorithmic
problems began to be formulated as problems of
existence of an algorithm for the solution of a given
infinite series of problems of a given type, and of find-
ing such an algorithm if it exists.
Several more exact definitions of algorithm appeared
in the theory of algorithms (cf. Algorithms, theory of)
practically all at the same time, and they all proved to
be essentially equivalent. In each such definition a suf-
ficiently broad class of specific algorithms, closed with
respect to the natural operation of combination of algo-
rithms, is singled out. Each statement to the effect that
some algorithmic problem is unsolvable is a precise and
proved mathematical theorem on the unsolvability of
the algorithmic problem under consideration by an
algorithm of the given class. In this form these
theorems may be regarded as specific, i.e. related to a
given class of algorithms. However, all results of this
kind may be expressed in terms of the intuitive idea of
an algorithm as understood by all mathematicians. This
'translation' is based on the so-called Church thesis
(depending on the exact specification of the given algo-

rithm, it may also be called Turing's thesis or the nor-
malization principle), which says that the class of algo-
rithms under consideration is universal, i.e. that the
essential function of any intuitive algorithm may be
performed by algorithms of this class. This thesis is a
scientific fact confirmed by the history of mathematics:
all algorithms known in mathematics satisfy it and all
attempts to construct algorithms not satisfying it have
failed. Finally, the thesis is also confirmed by the
equivalence of the various proposed definitions of algo-
rithm. Church' thesis cannot be proved, since it deals
with the mathematically vague concept of an intuitive
algorithm, but it is highly important nevertheless, since
it makes it possible to speak of unsolvability of given
algorithms in a wide sense, which is generally under-
stood by mathematicians.
The first examples of unsolvable algorithmic prob-
lems were noted in the theory of algorithms itself.
These include the problem of recognition of belonging
to a given non-recursive set of natural numbers, the
halting problem of a universal Turing machine, the
problem of the recognition of self-applicability of a
normal algorithm, etc.
As concerns algorithmic problems outside the theory
of algorithms itself, one must mention, first of all, the
problem of recognition of validity of formulas of first-
order predicate calculus, the unsolvability of which was
first proved by A. Church in 1936. Numerous studies of
algorithmic problems in model theory are related to
these results. Model theory, which deals with non-
empty sets with relations defined on them using predi-
cate calculus, was created in the nineteen-thirties by
A.I. Mal'tsev and A. Tarski. They must also be credited
with exact formulations of many algorithmic problems
in model theory and with several fundamental results in
this field. The elementary theory of a given class K of
models is the set of all closed formulas of first-order
predicate calculus that are true for all models of K. The
class K may consist of one model. An elementary
theory T is solvable or unsolvable, depending on whether
or not an algorithm for the recognition of whether or
not any given formula belongs to the theory T exists.
For a partial review of the methods of studying algo-
rithmic problems in model theory and the results
obtained, see [3]. The problem of the solvability of the
elementary theory of free groups is the most interesting
one mentioned in [3], and one which has not yet (1977)
been solved.
Numerous and varied algorithmic problems also arise
during the constructive interpretation of various
branches of mathematics. The principal results con-
cerning algorithmic problems in traditional branches of
mathematics are given below.
A naturally arising question remained open for a
ALGORITHMIC PROBLEM
long time: Are unsolvable algorithmic problems specific
to the theory of algorithms and mathematical logic? In
other words, do there exist unsolvable algorithmic
problems in traditional branches of mathematics far
removed from mathematical logic? The first result of
this kind was obtained in 1949 and 1947 by A.A. Mar-
kov and E.L. Post, independently of one another. They
proved that the problem of equality (identity) of words
for semi-groups - posed by A. Thue as early as 1914
- was unsolvable. In this problem one considers semi-
groups defined by a finite number of generators (an
alphabet)
(1)
and defining relations
A, = B;, i=l, ... ,k. (2)
An elementary transformation of the semi-group II here
considered is a transition from the word PAiQ to the
word PBiQ or vice versa, where P and Q are arbitrary
words in the alphabet (1). Two words X and Y in the
alphabet (1) are said to be equivalent in II if they coin-
cide, or if one can be obtained from the other through
a finite sequence of elementary transformations. The
set of all equivalence classes with a multiplication
operation which is naturally defined by the operation
of concatenation of words is just the semi-group II
defined by the generators (1) and relations (2). The
problem of word equality (identity) (the word problem) in
the semi-group II consists in finding an algorithm for
recognizing if any pair of words X and Y of II are
equivalent in II or not, i.e. if the elements of II defined
by them are identical or not. Markov and Post con-
structed specific semi-groups for which there is no algo-
rithm for solving the equality problem [1].
The most noteworthy result in this context was
obtained by P.S. Novikov [4], [5]. He showed that the
classical word problem in group theory (the equality or
identity of words problem) posed by M. Dehn in 1912,
which was studied by many experts in algebra
throughout the world, was unsolvable. It may be for-
mulated in the same way as the word problem for
semi-groups, except that only such systems of relations
(2) are considered for which inverse elements for all
generators (1) exist in the given semi-group. Novikov
also showed that the isomorphism problem, which is very
important in group theory - to wit, how to find an
algorithm to tell whether or not any two groups are iso-
morphic - is unsolvable. It was subsequently shown
that for any given group G it is impossible to find an
algorithm that would tell whether or not an arbitrary
group is isomorphic to G [6].
Markov studied the problem of recognition of invariant
properties of semi-groups, i.e. of properties that are
preserved on passing to isomorphic semi-groups [2]. He
145
ALGORITHMIC PROBLEM
showed that if for an invariant semi-group property a
there exists a semi-group with property a and another
semi-group which is not imbeddable in any semi-group
with this property, then there is no algorithm by which
it is possible to tell whether a given semi-group has the
property a or not. This actually means that almost-all
invariant semi-group properties are algorithmically
unrecognizable in the class of semi-groups. An analo-
gue of this fundamental theorem has also been
obtained in group theory [7], [8], [9]. In particular, the
conclusion in this case is as follows. Let Ka be the class
of all groups with an invariant property a. Two algo-
rithmic problems are related to each such class: the
word problem for groups from Ka, and the problem of
recognizing whether or not an arbitrary group belongs
to Ka. It was found that for any non-empty class Ka at
least one of these algorithmic problems is unsolvable.
The same applies to semi-groups as well.
Another problem studied was the specification of the
simplest groups and semi-groups in which the word
problem is unsolvable. Of the numerous studies on this
subject one may mention the proof of the unsolvability
of the word problem for the semi-group specified by
the following seven simple defining relations [10]:
ac = ca, ad = da, be = cb, bd = db,
eca = ce, edb = de, cca = ccae.
At a later date [11] an example was given of a semi-
group with three defining relations and an unsolvable
equality problem. No algorithm for solving the word
problem in the general case has yet (1977) been found,
not even for semi-groups with a single defining rela-
tion; it was only constructed for irreducible defining
relations [12]. The algorithm solving the word problem
for groups with one defining relation dates a long time
back [13], but even for two such relations the question
still remains open. The minimum number of defining
relations for which examples of groups with unsolvable
word problem are known is 12. The problem of word
equality for commutative semi-groups has been shown
to be solvable. For commutative groups the isomor-
phism problem has been solved as well. The problem
of word equality is solvable for finite groups, for
finitely-defined simple groups, for nilpotent groups and
for solvable groups of the second order. However, even
in the class of solvable groups of the fifth order one
can indicate a group given by a corresponding identity
of the fifth order and a finite number of defining rela-
tions with unsolvable word problem [15].
The study of the word problem for groups with a
limited measure of overlap of the defining words, i.e.
with overlap of the left-hand sides of the defining rela-
tions, had begun before the unsolvability of the word
problem in its general form was demonstrated. The
strongest result on this subject at the time of writing
146
(1977) is the proof of the solvability of the word prob-
lem in the class of groups defined by systems of defin-
ing relations in which the defining words may overlap
by less than 115 of their length [14]. If the overlap
between the words may be as large as 113 of their
length, an example of a group with unsolvable word
problem becomes quite simple. In the range between
115 and 1/3 the problem remains open. Similar studies
were also conducted for semi-groups. It was found that
for semi-groups with maximum measure of overlap of
the defining words less than 1/2 the word problem is
solvable, but examples are known of semi-groups with
an unsolvable word problem in which the maximum
overlap of the defining words is 1/2.
The most prominent topological algorithmic problem
is the classical decision problem for homeomorphism of
n-dimensional manifolds. Markov constructed an exam-
ple of a four-dimensional manifold we for which it is
impossible to construct an algorithm to tell whether or
not an arbitrary four-dimensional manifold is
homeomorphic to we [17]. It is known that for n = 2 the
solution of this problem is positive; for n = 3 the prob-
lem remains open; for four-dimensional and higher-
dimensional manifolds the problems of recognition of
combinatorial and homotopy equivalence of manifolds
are also unsolvable. One may also mention that if P is
a polyhedron for which the word identity problem of
its fundamental group is unsolvable, no algorithm can
be constructed for P that recognizes the bounded
homotopy (or the free homotopy) of two arbitrary
paths on P passing through a common point. A posi-
tive solution was obtained for the problem of conjuga-
tion in braid groups; this problem is equivalent to the
topological problem of recognizing the equivalence of
braids [16].
One of the most famous algorithmic problems in
mathematics is Hilbert's lO-th problem: To find an algo-
rithm by which to tell whether or not a system of
Diophantine equations with integer coefficients has a
solution in integers. After many examples of unsolv-
able algorithmic problems had become known, it was
suggested that this problem was unsolvable as well.
Moreover, a team of American mathematicians suc-
ceeded in proving that the un solvability of Hilbert's
tenth problem follows from the existence of a two-place
Diophantine predicate of exponential growth [18], [19].
However, they were unsuccessful in constructing such a
predicate; they merely demonstrated that the recogni-
tion of the existence of integer solutions of exponential
equations was an unsolvable problem. The desired
Diophantine predicate of exponential growth was first
constructed in 1970 [20], [21]; this was the first proof
ever given of the unsolvability of Hilbert's tenth prob-
lem.
In the theory of algorithms the existence of unsolv-

able algorithmic problems of various degrees of unsolva-
bility was proved. Many studies of algorithmic prob-
lems arising in mathematics showed that each such
problem had, in the general case, a maximum degree of
unsolvability, while their special instances (e.g. the
word identity problem in specific semi-groups or
groups) may have any pre-given degree of un solvability.
References
[I] MAL'TSEV, A.I.: AlgOrithms and recursive functions, Wolters-
Noordhoff, 1970 (translated from the Russian).
[2] MARKov, A.A.: Theory of algorithms, Israel Progr. Sci. Trans!.,
1961 (translated from the Russian). Also: Trudy Mat. Inst.
Steklov. 42 (1954).
[3] ERsHOV, Yu.L., Ef AL.: 'Elementary theories', Russian Math.
Surveys 20, no. 4 (1965), 35-105. (Uspekhi Mat. Nauk 20, no. 4
(1965),37-108)
[4] NOVIIWV, P.S.: Dokl. Akad. Nauk SSSR 85, no. 4 (1952),
709-712.
[5] NOVIKOV, P.S.: 'On algorithmical unsolvability of the word
problem', Trudy Mat. Inst. Akad. Nauk SSSR 44 (1955) (in
Russian).
[6] ADYAN, S.I.: 'Unsolvability of some algorithmic problems in
group theory', Trudy Moskov. Mat. Obshch. 6 (1957),231-298
(in Russian).
[7] ADYAN, S.I.: 'Finitely defined groups and automorphisms',
Dokl. Akad. Nauk SSSR 117, no. I (1957),9-12 (in Russian).
[8] ADYAN, S.l.: 'On algorithmic problems in effectively complete
classes of groups', Dokl. Akad. Nauk SSSR 123, no. I (1958),
13-16 (in Russian).
[9] RABIN, M.O.: 'Recursive unsolvability of group theoretic prob-
lems', Ann. of Math. (2) 67, no. 1(1958),172-194.
[IO] TSEiTIN, G.S.: 'An associative calculus with unsolvable
equivalence problem', Trudy Mat. Inst. Steklov. 52 (1958), 172-
189 (in Russian).
[II] MATIYASEVICH, YU.V.: 'Simple examples of undecidable asso-
ciative calculi', Soviet Math. Dokl. 8, no. 2 (1967), 555-557.
(Dokl. Akad. Nauk SSSR 173, no. 6 (1967), 1264-1266)
[I2] ADJAN, S.I. [S.I. ADYAN]: 'Defining relations and algorithmic
problems for groups and semigroups', Proc. Steklov Inst. Math.
85 (1967). (Trudy Mat. Inst. Steklov. 85 (1966
[I3] MAGNUS, W.: Uspekhi Mat. Nauk 8 (1941), 365-376.
[14] LYNDON, R.C.: 'On Dehn's algorithm', Math. Ann. 166 (1966),
208-228.
[15] REMESLENNIKOV, V.N.: 'Example of a finitely presented group
in the variety j5 with the unsolvable word problem', Alge-
bra and Logic 12, no. 5 (1973), 327-345. (Algebra i Logika 12,
no. 5 (1973), 577-602)
[16] GARSAID, F.A.: 'The braid group and other groups', Quart. J.
Math. Oxford Ser. (2) 20 (1969), 235-254.
[17] MARKov, A.A.: 'On the unsolvability of certain problems in
topology', Dokl. Akad. Nauk SSSR 123, no. 6 (1958), 978-980
(in Russian).
[18] DAVIS, M. and MASCHLER, M.: 'Reductions of Hilbert's tenth
problem', J. Symbolic Logic 23 (1958), 183-187.
[19] ROBINSON, 1.: 'Existential definability in arithmetic', Trans.
Amer. Math. Soc. 72 (1952), 437-449.
[20] MATIYASEVICH, YU.V.: 'Enumerable sets are Diophantine',
Soviet Math. Dokl. 11, no. 2 (1970), 354-357. (Dokl. Akad
Nauk SSSR 191, no. 2 (1970), 279-282)
[21] MATIYASEVICH, Yu.V.: 'Diophantine representation of enumer-
able predicates', Math. USSR-Izv. 5, no. I (1971), 1-28. (Izv.
Akad. Nauk SSSR Ser. Mat. 35 (1971),3-30)
S1. Adyan
Editorial comments. Church' thesis is often formulated
as stating that the idea of a recursive function (as defined
in various equivalent ways, e.g. via Turing machines, cf.
Turing machine), is the correct formalization of our intuitive
ALGORITHMIC REDUCIBILITY
concept of effective computability. Church' thesis is now
almost universally accepted.
The halting problem asks whether there exists an effective
procedure to decide, given a Turing machine M and input x,
whether the calculation ever terminates. A subset of Nk is
recursive if its characteristic function is recursive. Let all
Turing machines be enumerated and let Mx be the machine
with index x. Then the theorem that {(x, y): Mx with input y
halts} is not a recursive set is a precise statement of the
unsolvability of the halting problem.
The theorem to the effect that there exists a group with a
presentation (i.e. given by generators and relations) for
which the word problem is unsolvable is often known as the
Boone- Novikov theorem.
A class of groups @ is invariant if GE@ and H isomorphic
to G imply HE@. It is non-trivial if there exist finitely-
presented groups both inside and outside @ and it is
hereditary if a subgroup of an element of @ is an element of
@. The Adyan- Rabin theorem now says that if @ is an
invariant, non-trivial and hereditary class, then there is no
algorithm that determines whether the group given by a fin-
ite problem belongs to that class or not. Properties to which
the Adyan - Rabin theorem applies include being cyclic,
commutative, free, solvable, finite and equal to the identity
group.
The original papers in which Church' thesis was
advanced and defended are [A2] , [A3]. Together [A4] and
[A5] give an excellent introduction to algorithmic unsolvabil-
ity. Reference [A 1] contains a wealth of interesting material
and references on the word problem in groups.
References
[A1] BOONE, W.W., CANNONITO, F.B. and LYNDON, R.C.: Word
problems, North-Holland, 1973.
[A2] CHURCH, A.: 'An unsolvable problem of elementary number
theory (abstract)', Bull. Amer. Math. Soc. 41 (1935),332-
333.
[A3] CHURCH, A.: 'An unsolvable problem of elementary number
theory', Amer. J. Math. 58 (1936), 345-363.
[A4] DAVIS, M.: 'Unsolvable problems', in J. Barwise (ed.): Hand-
book of mathematical logic, North-Holland, 1977, pp. 567-
594.
[A5] ENDERTON, H.B.: 'Elements of recursion theory', in J. Bar-
wise (ed.): Handbook of mathematical logic, North-Holland,
1977, pp. 527-566.
AMS 1980 Subject Classification: 68020, 68025,
03DXX
ALGORITHMIC REDUCIBILITY - One of the basic
concepts in the theory of algorithms (cf. Algorithms,
theory of) and its applications. Its background was the
fact that solvability (and un solvability) of many algo-
rithmic problems (d. Algorithmic problem) is not esta-
blished directly, but rather by reducing some algo-
rithmic problem, already demonstrated to be unsolv-
able, to the problem in question, or by reducing the
latter problem to some other one which has already
been solved. Thus, the problem of homotopy of paths
in a polyhedron is shown to be unsolvable by reducing
the problem of equality of words in the corresponding
fundamental group to this problem.
147
ALGORITHMIC REDUCIBILITY
In what follows, the simplest examples of algorit1ullic
reducibility of number-theoretical predicates and func-
tions (i.e. those defined on natural numbers) will be
presented:
P(x) = Q(2x + 1),
where P and Q are predicates. The problem of solving
the predicate P, that is, of establishing the truth or
falsehood of P for different values of x, is reduced to
the problem of solving Q - or, 'P is reduced to Q', for
short. One may also say that the set of truths of P, i.e.
{x: P(x)}=Mp , is reduced to MQ={x: Q(x)}.
Let xl u
f(x) = ~ IIg(u,y),
u=ly=1
where f(x) and g(u,y) are number-theoretic functions.
The problem of computation of the function f is reduced
to the computation of g or, for short, f is reduced to g.
The concept of algorithmic reducibility was given a
more exact form by A.M. Turing: If, roughly speaking,
some Turing machine transforms a sequence of encoded
values of the function g into a sequence of encoded
values of the function f, then f is reduced to g. S.c.
Kleene [1] formulated the related concept of relative
computability with the aid of recursive systems of equa-
tions. Mter arithmetization, each algorithmic problem
is reduced to the problem of computing some number-
theoretic function f If f is reduced to g according to
Turing, f";;;;' Tg, and g is reduced to f, g,,;;, Tf, then one
says that f and g have the same degree of unsolvability
or that f Tg The relation ,,;;, T is both reflexive and
transitive. Thus, all functions (and sets of natural
numbers or their characteristic predicates) are subdi-
vided into equivalence classes, known as Turing degrees
or T-degrees [3]. Most algorithmic problems considered
in logic and in mathematics are expressible as problems
of membership in enumerable sets of constructive
objects. In this context the study of enumerable sets
was initiated in the nineteen-fourties by E.L. Post [2];
he introduced certain special kinds of algorithmic redu-
cibility, in addition to Turing's kind, and formulated
the problem of reducibility as follows: Can various enu-
merable unsolvable sets be reduced (according to Tur-
ing) to each other? It was subsequently established that
the enumerable sets form an infinite, very rich system
of T-degrees. This was supplemented by discovering the
so-called priority method, which is extensively employed
in the theory of algorithms.
Subsequently, degrees of unsolvability found applica-
tions in other branches of mathematics as well. For
instance, for all T-degrees a and b of enumerable sets,
if a";;' Tb, then there exists a finitely-determined group
in which the problem of equality of words has degree a,
while the conjugation problem has degree b. There is a
close connection between degrees (concerning the dif-
148
ferent kinds of algorithmic reducibility) of enumerable
sets and the rate of growth of the complexity of initial
fragments of enumerable sets. A special kind of reduci-
bility - polynomial reducibility or p-reducibility (reduci-
bility with some limitation imposed on time) - is
employed in proving the universal character of com-
binatorial optimization problems from various branches
of mathematics [5]; for more details on this subject see
[1]. The methods of measure theory [1], and the forcing
method [4] have in turn been used in the study of
degrees.
References
[I) ROGERS, JR., H.: Theory of recursive functions and effective
computability, McGraw-Hill, 1967.
[2) POST, E.L.: 'Recursively enumerable sets of positive integers
and their decision problems', Bull. Amer. Math. Soc. 50 (1944),
284-316.
[3) KLEENE, S.c. and POST, E.L.: 'The upper semi-lattice of
degrees of recursive unsolvability', Ann. of Math. (2) 59, no. 3
(1954), 379-407.
[4) SELMAN, A.L.: 'Applications of forcing to the degree-theory of
the arithmetical hierarchy', Proc. London Math. Soc. (3) 2S
(1972), 586-602.
[5) KARP, R.M.: 'Reducibility among combinatorial problems', in
R.E. Miller and l.W. Thatcher (eds.): Complexity of computer
computations, Plenum Press, 1972, pp. 85-103. A A M h 'k
.. uc m
Editorial comments.
References
[A 1] R.I.: Recursively enumerable sets and degrees,
SOARE,
Springer, 1986.
[A2] SIMPSON, S.G.: 'Degrees of unsolvability: a survey of results',
in J. Barwise (ed.): Handbook of mathematical logic, North-
Holland, 1977, pp. 631-652.
AMS 1980 Subject Classification: 03DXX, 03FXX,
68CXX
ALGORITHMIC THEORY OF SETS, algorithmic set
theory - See Recursive set theory.
AMS 1980 Subject Classification: 03E45
ALGORITHMS, COMBINATIONS OF - The name of
several methods for constructing new algorithms (cf.
Algorithm) from given algorithms.
As applied to normal algorithms (cf. Normal algo-
rithm), the following combinations (compositions) are
best known: normal composition (IB o 2r) of two normal
algorithms 2r and IB, normal union (2r" IB) of two nor-
mal algorithms 2r and IB, normal branching (2rYIBI G) of
two normal algorithms 2r and IB controlled by a normal
algorithm G, and normal repetition 9! OG of a normal
algorithm 9! controlled by a normal algorithm G. If
9!, IB and G are normal algorithms in some alphabet A,
the above combinations are normal algorithms in a cer-
tain given extension of A and satisfy the following con-
ditions: a) for any word P in A, IBo9!)lP J~I.B l2r lP JJ
is true (the composition theorem); b) for any word P in
A, ( 2r "IS lP J~2r lP JIBlP J is true (the union theorem);

c) for any word P in A
(ID'IS IC) lP J ~ { ill:lPJ
IS lP J
ifClPJoF,
if C lP J$:F,
moreover, if (~?B I~ lP J has been defined, then ~ lP J
is defined as well (the branching theorem); d) for any
words P and Q in A, the graphic equality
(~O~HP J Q is true if and only if it is possible to
indicate a sequence of words Po, ... ,Pk (k ~ I) over
the alphabet A such that
Po P, Pk ::g:: Q,
Pi ~ ill:lPi-d (i=l, ... ,k),
ClPkJ $: F (i=l, ... ,k-l),
ClPkJ ::g:: F
(the repetition theorem). Similar theorems may also be
obtained for Turing machines (cf. Turing machine). In
the theory of recursive functions (cf. Recursive func-
tion), the combinations of these functions which are
most frequently employed are those supplied by the
substitution operator, the primitive recursion operator
and the /L-operator.
Theorems on combinations of algorithms reveal an
important feature of the existing formalizations of the
general concept of an algorithm - to wit, their 'clo-
sure' under the natural ways of combination of algo-
rithms. This fact is one of the principal arguments in
favour of the basic assumption about algorithms (the
Church thesis). Theorems on compositions of algo-
rithms are an important part of the general theory of
algorithms. Having been demonstrated once, they make
it possible to determine the realizability of complicated
and cumbersome algorithms, without actually writing
down their schemes.
Of major interest in the general theory of algorithms
is the problem of synthesizing an arbitrary algorithm
within some class of interest, using a pre-determined set
of composition operators.
References
[1] MARKov, A.A.: Theory of algorithms, Israel Progr. Sci. Trans!.,
Jerusalem, 1961 (translated from the Russian). Also: Trudy
Mat. Inst. Steklov. 42 (1954).
[2] KLEENE, S.c.: Introduction to metamathematics, North-Holland,
1951.
[3] OUSPENSKl, V.A. [VA USPENSKlIl: Le~ons sur les fonctions cal-
culable, Hermann, 1966 (translated from the Russian).
[4] MAL'TSEV, A.I.: Algorithms and recursive functions, Wolters-
Noordhoff, 1970 (translated from the Russian).
N.M. Nagornyr
Editorial comments. The standard terminology for the
above composition operators is in the West: 'sequential
composition', 'non-deterministic choice', 'if-then-else' and
the 'while loop', in order of appearance. A representative
result of the type indicated is the characterization of the
power of structured programming: Every flowchart can be
expressed as a flowchart obtained by sequential composi-
ALGORITHMS. EQUIVALENCE OF
tion, if-then-else and while, from simple assignment state-
ments [A1].
References
[A1] BOHM, C. and JACOPINI, G.: 'Flow diagrams, Turing machines
and languages with only two formation rules', Comm. ACM 9
(1966), 366-371.
AMS 1980 Subject Classification: 68C05
ALGORITHMS, EQUIVALENCE OF - A binary rela-
tion connecting algorithms of a fixed type and express-
ing the fact that any two algorithms thus connected
yield the same results if fed with the same given type of
inputs (and may also yield some additional information
as regards the computations thus performed - the so-
called history of the computation). A few typical exam-
ples of this connection are given below.
a) Consider all possible recursive schemes - i.e. sys-
tems of equations which define n-place partial recursive
functions. Two schemes which define functions <I> and
1/; are called functionally eqUivalent if for any natural
numbers Xl, . . . ,Xn the conditional equality
<P(X1,'" ,xn ) ~ 1/-{X 1, . . ,xn )
is true (it is considered to be true if both its sides are
meaningful at the same time and, if meaningful, assume
the same value). It was shown by S.c. Kleene (see, for
example, [ID that for any everywhere-defined comput-
able operator lR over recursive schemes it is possible to
find a scheme S such that S and lR(S) are functionally
equivalent (the so-called fixed-point theorem). The
consequence of this is, in particular, the
Uspenskii - Rice theorem on the un solvability of any
property of recursive schemes which is invariant with
respect to functional equivalence if there exist schemes
S 1 and S 2 such that S 1 has this property, while
scheme S 2 does not have it. The un solvability of the
relation of functional equivalence itself follows from
this theorem.
b) Consider normal algorithms (cf. Normal algorithm)
over a fixed alphabet A. Two such algorithms ~ and ?B
are called equivalent with respect to A [2] if for any
word P in A the following condition is satisfied: If one
of these algorithms converts P to some word Q which is
also contained in the alphabet A, then the second algo-
rithm also converts P to Q. These algorithms are called
completely equivalent with respect to A if for any P in A
the conditional equality
ill:(P) ~ IS(P)
is true. Both these relations are unsolvable.
c) Consider logical schemes of algorithnls (Yanov
schemes, [3D. The realization of such a scheme is a
sequence of operators involved in its execution, from
beginning to end. Two schemes are called equivalent if
the sets of their realizations coincide. The equivalence
149
ALGORITHMS, EQUIVALENCE OF
relation of Yanov schemes proved to be solvable, and a
complete system of equivalent transformations has been
constructed for it [3], [4].
A detailed study of the relations of algorithmic
equivalence is of major importance in the context of a
number of tasks of current interest (especially in the
theory of programming schemes), which can only be
solved by advanced techniques of equivalent transfor-
mations of algorithms. Such problems include the
translation of algorithms from one language to another
and their optimization (transformations aimed at
improving their 'working characteristics'). Special atten-
tion is given in this connection [5] to finding solvable
relations of algorithmic equivalence which would per-
mit complete systems of equivalent transformations
that are convienient with regard to their applications.
The concept outlined in [3] to a large extent determined
the general approach to the study of algorithmic
equivalence.
References
[I] KLEENE, S.c.: Introduction to metamathematics, North-Holland,
I95l.
[2] MARKov, A.A.: Theory of algorithms, Israel Progr. Sci. Trans!.,
1961 (translated from the Russian). Also: Trudy Mat. Inst.
Steklov. 42 (1954).
[3] YANOV, Yu.!.: 'On logical schemes of algorithms', Problemy
Kibernet. 1 (1958), 75-127 (in Russian).
[4] ERSHOV, A.P.: 'Operator algorithms III. (On Yanov operator
schemes)" Problemy Kibernet. 20 (1968), 181-200 (in Russian).
[5] ERSHOV, A.P.: 'The contemporary state of the theory of pro-
gramming schemes', Problemy Kibernet. 27 (1973),87-110 (in
Russian). NM N
. . agornYI
AMS 1980 Subject Classification: 68C05
ALGORITHMS, THEORY OF - The branch of
mathematics dealing with the general properties of
algorithms (cf. Algorithm). Preliminary notions about
algorithms have played a role in mathematics
throughout history. However, the concept of an algo-
rithm itself was only formulated in the twentieth cen-
tury and became the object of independent study (at
first only in its rather vaguely defined form) in the
nineteen-twenties by the intuitionistic school of L.E.J.
Brouwer and H. Weyl [1] (cf. Intuitionism). Systematic
studies began in 1936, with the publication by A.
Church [2] of the first formalization of the concept of a
computable function (he suggested that the concept of
an everywhere-defined computable function with
natural arguments and natural values be identified with
the concept of a general recursive function) and of the
first example of a non-computable function. A.M. Tur-
ing [3], [4] and E.L. Post [5] gave the first formaliza-
tions of the concept of an algorithm (in terms of ideal-
ized computers, cf. Turing machine). Subsequent
development of the theory of algorithms is due to the
studies of Kleene, Post [6], [7], [8], A.A. Markov [9],
150
[10], [11], and others. In particular, Markov rendered
the concept of an algorithm more precise by the intro-
duction of the concept of a normal algorithm [10]. A.N.
Kolmogorov [12], [13] suggested the most general
approach to the concept.
Since algorithms exclusively deal with the so-called
constructive objects (cf. Constructive object), the ideas
and methods of the theory of algorithms can be applied
to non-constructive objects only if these are encoded,
or re-named, as constructive objects. The study of the
general properties of such encodings (mainly of situa-
tions in which such encodings, or names, are natural
numbers) is the subject of the theory of enumeration
[14], which is an important part of the theory of algo-
rithms.
Fundamental concepts in the theory of algorithms. The
domain of applicability of an algorithm is the set of the
objects to which it is applicable. It is said that an algo-
rithm ~ 1) 'computes a function l' if its domain coin-
cides with the range of definition of f, and if ~ con-
verts any x from its domain into f (x); 2) 'solves a set
A with respect to a set X' if it is applicable to any x
from X and converts all x from X n A to the word
'yes', and all x from X\A to the word 'no'; and 3)
'counts (enumerates) a set B' if its domain is the
sequence of natural numbers, and the set of results is
B. A function is called computable if an algorithm
which evaluates the function exists. A set is said to be
solvable with respect to X if there exists an algorithm
-
which solves it with respect to X. A set is said to be
enumerable if it is empty or if there exists an algorithm
enumerating it.
A detailed analysis of the concept of an algorithm
reveals that: 1) the range of possible initial data and
the range of applicability of an algorithm are enumer-
able sets. In tum, 2) for any pair of enumerable sets,
one of which is included in the other, it is possible to
find an algorithm in which the larger set is the range of
possible inputs, while the smaller set is the range of
outputs. The following basic theorems hold: 3) a func-
tion f is computable if and only if its graph, i.e. the set
of all pairs of the form <x, f(x, is enumerable; 4) a
subset A of an enumerable set X is solvable with
respect to X if and only if A and X \ A are enumer-
able; 5) if A and B are enumerable, A U B and A n B
are also enumerable; 6) in each infinite enumerable set
X there exists an enumerable subset with a non-
enumerable complement (in view of 4) this enumerable
subset will not be solvable with respect to X); and 7)
for each infinite enumerable set X there exists a com-
putable function, defined on a subset of this set, which
cannot be extended to a computable function defined
on all of X. Statements 6) and 2) in combination yield
an example (d. Algorithm) of an algorithm with an

unsolvable domain of applicability. Solvable and enu-
merable sets represent the simplest (and the most
important) examples of sets, the structure of which is
defined by certain algorithmic procedures. A systematic
study of sets of constructive objects from the point of
view of the properties of these sets as related to the
availability of some of algorithms, constitutes the so-
called algorithmic theory of sets; certain concepts,
methods and results of this theory find analogues in
descriptive set theory.
Algorithmic problems. The problem of constructing an
algorithm with certain given properties is known as the
algorithmic problem. As a rule, the property of the
algorithm which is sought is formulated in terms of the
correspondence which should hold between the inputs
and the results of the algorithm. Some important
examples of algorithmic problems are: the problem of
computing a given function (i.e. finding an algorithm
which computes the function); the problem of solving a
given set (i.e. finding an algorithm which solves this
set with respect to some other set); and the problem of
enumeration of a given set (i.e. to find an algorithm
enumerating the given set). All the examples of algo-
rithmic problems in various branches of mathematics,
given in the section 'Applications' below, are problems
of solving. That an algorithmic problem is unsolvable
means that there is no corresponding algorithm;
theorems which establish the un solvability of such
problems belong to the most important theorems in the
theory of algorithms. For instance, if an algorithm has
an unsolvable domain, the algorithmic problem of solv-
ing this domain with respect to the set of all possible
inputs is unsolvable. By reducing to this problem unsol-
vability, most other solvability problems were found to
be unsolvable; this applies, in particular, to all the
problems listed in the section 'Applications' below. The
question as to whether or not the unsolvability of any
unsolvable problem can be established in this way,
forms the so-called algorithmic reducibility problem.
The metric theory of algorithms. The theory of algo-
rithms comprises the descriptive (qualitative) and the
metric (quantitative) theory. The former deals with
algorithms from the point of view of the correspon-
dence they establish between the initial data and the
results; in particular, all algorithmic problems discussed
in the preceding section belong to the qualitative
theory. The latter studies algorithms from the point of
view of the complexity both of the algorithms them-
selves (cf. Algorithm, complexity of description of an)
and of the 'computations' defined by the algorithms,
i.e. the processes of successive transformations of con-
structive objects (cf. Algorithm, computational complex-
ity of an). It is important to note that both the compu-
tational complexity of the algorithm and the complexity
ALGORITHMS, THEORY OF
of description may be defined in various ways, and it
may well be that A will tum out to be more complex
than B when one definition is adopted, while B will
tum out to be more complex than A according to
another definition. In order to be able to speak about
the complexity of algorithms, one must first specify
some exact language in which the algorithms are to be
written, and define the complexity of an algorithm as
the complexity of its notation; the complexity of nota-
tion can in tum be determined in various ways (e.g. as
the number of symbols of a given type which are
included in the notation, or as a selection of such
numbers computed for the various types of symbols).
In order to speak about computational complexity, one
must specify the exact form of the computation as a
chain of constructive objects successively replacing each
other, as well as some criterion of the complexity of
such a chain - the number of 'links' (or steps) in the
chain alone or in combination with the 'dimension' of
the links, etc. In any event, the computational complex-
ity will depend on the input with which the computa-
tion is begun; for this reason, the computational com-
plexity is a function which assigns to each object within
the domain of the algorithm the complexity of the
corresponding chain. The development of the methods
of evaluation of the complexity of algorithms and of
computations has a high theoretical and practical
importance; however, as distinct from the descriptive
theory of algorithms, which has now crystallized into
an integral mathematical discipline [11], [15], [16], the
metric theory of algorithms is only in the process of
being created [17], [18], [19], [20].
Applications of the theory of algorithms. These are
encountered in all branches of mathematics in which
algorithmic problems are encountered. Such problems
may arise in mathematical logic and in model theory;
for each theory, the problem which arises is to solve the
set of all true or demonstrable statements of this theory
with respect to the set of all its statements (the theories
divide into solvable or unsolvable theories, depending
on the solvability or unsolvability of the problem in
question). Church [21], [22] proved in 1936 that the
problem of solving the set of all true assumptions of
predicate logic was unsolvable; other important results
on the subject are due to A. Tarski, A.I. Mal'tsev and
to others [23]. Unsolvable algorithmic problems are
encountered in algebra (the word (identity) problem for
semi-groups and, in particular, for groups). The first
examples of semi-groups with unsolvable word problem
were found, independently of each other, by Markov [9]
and by Post [8] in 1947, while an example of a group
with unsolvable word problem was found by P.S. Novi-
kov [24], [25] in 1952. Markov [26] in 1958 proved that
the problem of homeomorphy in topology was unsolv-
151
ALGORITHMS, THEORY OF
able for an important class of cases. In number theory,
Yu.V. Matiyasevich proved in 1970 that the problem of
solvability of Diophantine equations was unsolvable
[27], [28]. Similar examples can also be quoted from
many other branches of mathematics.
The theory of algorithms is closely connected with
mathematical logic, since the concept of an algorithm
forms the base of one of the central concepts of
mathematical logic - the concept of a calculus, as a
result of which the GOdel incompleteness theorem of
formal systems may be obtained from theorems of the
theory of algorithms [29]. Finally, the theory of algo-
rithms is closely connected with the foundations of
mathematics, where one of the key problems is the rela-
tion between the constructive and the non-constructive.
In particular, the theory of algorithms provides the
apparatus for the development of the constructive
direction in mathematics. It was suggested by KoImo-
gorov [30] in 1965 that the theory of algorithms be used
as the foundation of information theory (cf. Algo-
rithmic information theory). The theory of algorithms is
the theoretical foundation for a number of problems in
computational mathematics, and is closely related to
cybernetics, in which an important subject is the study
of control algorithms.
References
[I] WEYL, H.: Philosophie der Mathematik und NaturwissenschaJ-
ten, Leibniz Verlag, Miinchen, 1950.
[2] CHURCH, A.: 'An unsolvable problem of elementary number
theory', Amer. J. Math. 58 (1936),345-363.
[3] TuRING, A.M.: 'On computable numbers, with an application
to the Entscheidungsproblem', Proc. London Math. Soc. (2) 42
(1937), 230-265.
[4] TuRING, A.M.: 'On computable numbers with an application
to the Entscheidungsproblem, a correction', Proc. London
Math. Soc. (2) 43 (1937),544-546.
[5] POST, E.L.: 'Finite combinatory processes - formulation I', J.
Symbolic Logic 1, no. 3 (1936), 103-105.
[6] POST, E.L.: 'Formal reductions of the general combinatorial
decision problem', Amer. J. Math. 65 (1943), 197-215.
[7] POST, E.L.: 'Recursively enumerable sets of positive integers
and their decision problems', Bull. A mer. Math. Soc. 50 (1944),
284-316.
[8] POST, E.L.: 'Recursive unsolvability of a problem of Thue', J.
Symbolic Logic 12, no. I (1947), I-II.
[9] MARKov, A.A.: Dokl. Akad. Nauk SSSR 55, no. 7 (1947),
587-590.
[10] MARKov, A.A.: 'Theory of algorithms', Trudy Mat. Inst. Stek-
lov. 38 (1951), 176-189 (in Russian).
[11] MARKov, A.A.: Theory of algorithms, Israel Progr. Sci. Trans!.,
1961 (translated from the Russian). Also: Trudy Mat. Inst.
Steklov. 42 (1954).
[12] KOLMOGOROV, A.N.: 'On the concept of an algorithm',
Uspekhi Mat. Nauk 8, no. 4 (1953), 175-176 (in Russian).
[13] KOLMOGOROV, A.N. and USPENSKIi, V.A.: 'On the definition
of an algorithm', Uspekhi Mat. Nauk 13, no. 4 (1958),3-28 (in
Russian).
[14] ERSHOV, YU.L.: The theory of enumeration. 1-2 .. Novosibirssk,
1969-1973 (in Russian).
[15] MAL'TSEV, A.I.: Algorithms and recursive functions, Wolters-
Noordhoff, 1970 (translated from the Russian).
[16] ROGERS, JR., H.: Theory of recursive functions and effective
computability, McGraw-Hili. 1967.
152
[17] MARKov, A.A.: 'Normal algorithms connected with the com-
putation of Boolean functions', Izv. Akad. Nauk SSSR Ser.
Mat. 31 (1967), 161-208 (in Russian).
[18] 'fRAKHTENBROT, B.A.: Complexity of algorithms and calcula-
tions, Novosibirssk, 1967 (in Russian).
[19] Problems of mathematical lOgic. Complexity of algorithms and of
classes of computable functions, Moscow, 1970 (in Russian;
translated from the English).
[20] Complexity of computation of algorithms, Moscow, 1974 (in
Russian; translated from the English).
[21] CHURCH, A.: 'A note on the Entscheidungsproblem', J. Sym-
bolic Logie 1, no. I (1936), 40-41.
[22] CHURCH, A.: 'Correction to 'a note on the Entscheidungsprob-
lem", J. Symbolic Logic 1, no. 3 (1936), 101-102.
[23] ERsHOV, Yu.L., ET AL.: 'Elementary theories', Russian Math.
Surveys 20, no. 4 (1965),35-105. (Uspekhi Mat. Nauk 20, no. 4
(1965),37-108)
[24] NOVIKOV, P.S.: 'On algorithmic unsolvability of the word
problem', Dokl. Akad. Nauk SSSR 85, no. 4 (1952),709-712
(in Russian).
[25] NOVIKOV, P.S.: On algorithmically unsolvable problems of word
identity in groups, Moscow, 1955 (in Russian).
[26] MARKov, A.A.: 'Unsolvability of the homeomorphy problem',
Dokl. Akad. Nauk SSSR 121, no. 2 (1958), 218-220 (in Rus-
sian).
[27] MATIYASEVICH, Yu.V.: 'The Diophantiness of enumerable
sets', Dokl. Akad. Nauk SSSR 191, no. 2 (1970), 279-282 (in
Russian).
[28] MATIYASEVICH, YU.V.: 'Diophantine sets', Russian Math. Sur-
veys 27, no. 5 (1972), 124-164. (Uspekhi Mat. Nauk 27, no. 5
(1972), 185-222)
[29] USPENSKIi, V.A.: 'An elementary exposition of GOdel's incom-
pleteness theorem', Russian Math. Surveys 29, no. I (1974),63-
106. (Uspekhi Mat. Nauk 29, no. I (1974), 3-47)
[30] KOLMOGOROV, A.N.: 'Three approaches to the definition of
the concept of 'quantity of information", Problemy Peredachi
Informatsii 1, no. 1 (1965), 3-11 (in Russian). VA U. k
.. spens II
Editorial comments. Since the above text was prepared
also the theory of complexity of algorithms has matured into
a well-developed branch of mathematics.
See also the comments to Algorithm, computational
complexity of an.
In addition to [27], [28] one may consult [A1J, [A2]. A
thorough treatise on computational complexity is [A3].
References
[A1] DAVIS, M.: 'Hilbert's tenth problem is unsolvable', Amer.
Math. Monthly 80 (1973), 233-269.
[A2] MAmAsEvIc, J.V.: 'A new proof of the theorem of exponen-
tial Diophantine representation of enumerable sets', J. Soviet
Math. 14 (1980),1475-1486.
[A3] WAGNER, K. and WECHSUNG, G.: Computational complexity,
Reidel, 1986.
AMS 1980 Subject Classification: 68CXX, 03DXX
ALIQUOT RATIO, aliquot fraction - A fraction of
the type I / n, where n is a natural number. It is essen-
tial, in solving several physical and mathematical prob-
lems, that each positive rational number is represent-
able as a sum of a finite number of aliquot ratios, with
different denominators. Thus. 3/ II = 1/6 + 1/11 +
1/66. Aliquot ratios were extensively employed III
Ancient Egypt and, as a result, received the name of
Egyptian fractions.
B. M. Bredikhin
AMS 1980 Subject Classification: 10A99

ALMOST-COMPLEX STRUCTURE - A tensor field
I of linear transformations of the tangent spaces on a
manifold M satisfying the condition
12 = -id,
i.e. a field of complex structures in the tangent spaces
TpM, p EM. An almost-complex structure I dete~es
a decomposition T C M = V + + V_of the complexiflca-
tion rc M of the tangent bundle in a direct sum of two
complex mutually-conjugate subbundles V + and V_
consisting of eigen vectors of the affinor I (extended by
linearity on rc M) with eigen values i and - i, respec-
tively. Conversely, a decomposition of TC M in a direc.,!
sum of mutually-conjugate vector subbundles S, S
defines an almost-complex structure on M for which
V+=S.
An almost-complex structure I is called integrable if
it is induced by a complex structure on M, i.e. if there
exists an atlas of admissible charts of the manifold M
in which the field I has constant coordinates Ii. A
necessary and sufficient condition for the integrability
of an almost-complex structure is that the subbundle
V + is involutive, i.e. that the space of its sections is
closed with respect to commutation of (complex) vector
fields. The condition for the subbundle V + to be invo-
lutive is equivalent to the vanishing of the vector-
valued 2-form N(I, I) associated with I and given by
the formula
N(l, lXX, Y) = [IX, Iy]-I[X, IY]-I[l, XYj-[X, y],
where X and Y vector fields. This form is called the tor-
sion tensor, or the Nijenhuis tensor, of the almost-
complex structure. The torsion tensor N(I, I) can be
considered as first-order differentiation on the algebra
of differential forms on M of the form
N(I, I) = [1,[1, dll+d,
where d is the exterior differential and I is considered
as a differentiation of order zero.
From the point of view of the theory of G-structures
an almost-complex structure is a GL(m, C)-structure,
where m = (1 / 2)dim M, and the torsion tensor N (I, I)
is the tensor defined by the first structure function of
this structure. A GL(m, C)-structure is a structure of
elliptic type, therefore the Lie algebra of infinitesimal
automorphisms of an almost-complex structure satisfies
a second-order system of elliptic differential equations
[1]. In particular, the Lie algebra of infinitesimal auto-
morphisms of an almost-complex structure on a com-
pact manifold is finite-dimensional, and the group G of
all automorphisms of a compact manifold with an
almost-complex structure is a Lie group. For non-
compact manifolds this statement is, in general, not
true.
If the automorphism group G acts transitively on the
manifold M, then the almost-complex structure I is
ALMOST-PERIOD
uniquely defined by its value 1p at a fixed point p EM.
This represents an invariant complex structure in the
tangent space TpM relative to the isotropic representa-
tion (see Invariant object on a homogeneous space).
Methods of the theory of Lie groups allow one to con-
struct a wide class of homogeneous spaces having an
invariant almost-complex structure (both integrable and
non-integrable) and to classify invariant almost-
complex structures under different assumptions (see
[2]). For instance, any quotient space G / H of a Lie
group G by the subgroup H consisting of fixed points
of an automorphism of even order of G has an invari-
ant almost-complex structure. An example is the 6-
dimensional sphere S6, considered as the homogeneous
space G2 / SU(3); none of the invariant almost-
complex structures on S6 is integrable.
The existence of an almost-complex structure on a
manifold imposes certain restrictions on its topology -
it must be of even dimension, oriented, and in the com-
pact case all its odd-dimensional Stiefel- Whitney
classes must vanish. Among the spheres only the
spheres of dimensions 2 and 6 admit an almost-
complex structure.
References
[1] KOBAYASHI, S.: TransfoT7lUJtion groups in differential geometry,
Springer, 1972.
[2] KOMRAKOY, B.P.: Structure on manifolds and homogeneous
spaces, Minsk, 1978 (in Russian).
[3] LICHNEROWICZ, A.: Global theory of connections and holonomy
groups, Noordhoff, 1976 (translated from the French).
[4] WELLS, JR., R.O.: Differential analysis on complex manifolds,
Springer, 1980.
[5] HORMANDER, L.: An introduction to complex analysis in several
variables, North-Holland, 1973. D. V. Alekseevskii
Editorial comments. The theorem that an almost-
complex structure is integrable, i.e. comes from a complex
structure, if and only if its Nijenhuis tensor vanishes, is due
to A. Newlander and L. Nirenberg [A1].
References
[A1] NEWLANDER, A. and NIRENBERG, L.: 'Complex analytic coor-
dinates in almost complex manifolds', Ann. of Math. 65
(1957), 391-404.
AMS 1980 Subject Classification: 53C15, 58812
ALMOST-EVERYWHERE, for almost all x (with
respect to a measure JL) - An expression indicating that
one speaks about all x of a measure space X with the
possible exception of some set A ~X of measure zero:
JL{A)=O. V.I. Bityutskov
AMS 1980 Subject Classification: 28-XX
ALMOST-PERIOD - A concept from the theory of
almost-periodic functions (cf. Almost-periodic func-
tion); a generalization of the notion of a period. For a
uniformly almost-periodic function f(x), - oo<x<oo,
a number T=Tff.) is called an f.-almost-period of f(x) if
153
ALMOST-PERIOD

for all x, almost-periodic analytic functions turns out to be a

If(X +r)-f(x) I < (.
For generalized almost-periodic functions the concept
of an almost-period is more complicated. For example,
in the space SIf an (-almost-period T is defined by the
inequality D.1,[f(x+'T),f(x)) < (,
where DSj [f, cf>] is the distance between f(x) and cf>(x) in
the metric of SIf.
A set of almost-periods of a function f(x) is said to
be relatively dense if there is a number L=L(,jO
such that every interval (a, a+ L) of the real line con-
tains at least one number from this set. The concepts of
uniformly almost-periodic functions and that of
Stepanov almost-periodic functions may be defined by
requiring the existence of relatively-dense sets of (-
almost-periods for these functions.
References
[I) LEVITAN, B.M.: Almost-periodic functions, Moscow, 1953 (in
Russian).
E.A. Bredikhina
Editorial comments. For the definition of Sf and its
metric Os: see Almost-periodic function. The Weyl,
Besicovitch and Levitan almost-periodic functions can also
be characterized in terms of Sf (-periods. These characteri-
zations are more complicated. A good additional reference
is [A1), especially Chapt. II.
References
[A1) BESICOVITCH, A.S.: Almost periodic functions, Cambridge
Univ. Press, 1932.
AMS 1980 Subject Classification: 42A75, 43A60
ALMOST-PERIODIC ANALYTIC FUNCTION - An
analytic function f(s), S=(J+iT, regular in a strip
- 00 ';;;;'a<(J<f3.;;;;, + 00, and expandable into a series
where the an are complex and the An are real numbers.
A real number T is called an (-almost-period of f(s) if
for all points of the strip (a,f3) the inequality
If(s+i'T)-f(s) I <
holds. An almost-periodic analytic function is an ana-
lytic function that is regular in a strip (a,f3) and
possesses a relatively-dense set of (-almost-periods for
every (>0. An almost-periodic analytic function on a
theory analogous to that of almost-periodic functions of
a real variable (cf. Almost-periodic function). Therefore,
many important results of the latter theory can be
easily carried over to almost-periodic analytic func-
tions: the uniqueness theorem, Parseval's equality, rules
of operation with Dirichlet series, the approximation
theorem, and several other theorems.
References
[1) BOHR, H.: Almost-periodicfunctions, Chelsea, reprint, 1947
(translated from the German).
[2) LEVITAN, B.M.: Almost-periodic functions, Moscow, 1953,
Chapt. 7 (in Russian).
E.A. Bredikhina
Editorial comments. The hyphen between almost and
periodic is sometimes dropped.
References
[A1] CORDUNEANU, c.: Almost periodic functions, Interscience,
1961, Chap!. 3.
AMS 1980 Subject Classification: 42A75, 43A60,
30H05, 30EXX
ALMOST-PERIODIC FUNCTION - A function
representable as a generalized Fourier series. There are
several ways of defining classes of almost-periodic func-
tions, based respectively on notions of closure, of an
almost-period and of translation. Each of these classes
can be obtained as a closure, with respect to some
metric, of the set of all finite trigonometric sums.
Let DG[f(x), cf>(x)] be the distance of two functions
f(x) and cf>(x) in a metric space G of real- or complex-
valued functions on R. In the following, G will be one
of the spaces U, S, WP, or BP. Here U is the set of
continuous bounded functions on the real line with the
metric
Du[f(x), <j>(x) = sup
-oo<x<oo
If(x)-<j>(x) I;
and S, WP and BP for p;;:.1 are the sets of functions
that are measurable and whose p-th powers are integr-
able on every finite interval of the real line, the metrics
being
Ds~[f(x), <j>(x) =
IX+/ lilP
_oos~p<oo [ I [ I f(x)-<j>(x) IP dx ,
Dwp[f(x, <j>(x)] = lim D\~[f(x). <j>(x)],
!----x

closed strip a';;;;'(J.;;;;,f3 is defined similarly. An almost-

periodic analytic function on a strip [lX, 13] is a uni-
Dsp[f(x), <t>(x)] = [

T_X
1
lim -2
7T
f I j(x)-<j>(x) IP dxlilP.
T

--7

formly almost-periodic function of the real variable 'T

Let T be the set of trigonometric polynomials
on every straight line in the strip and it is bounded in
v
[lX, 13], i.e. on any interior strip. If a function f(s), regu- ~
L." ak e
IA~ ,-
.
lar in a strip (a, 13), is a uniformly almost-periodic func- k - I

tion on at least one line 0 = 00 in the strip, then bound- where the Ak are arbitrary real numbers and the Uk are
edness of J(s) in [a,f3] implies its almost-periodicity on complex coefficients, and let the symbol Hc(T) denote
the entire strip [a, /3]. Consequently, the theory of the closure of T in G. The classes H veT) = U-a.p.,

154

HSj(T)=SI/-a.p., HW"(T) = wP-a.p. and HRP = BP-a.p.
denote, respectively, the classes of uniformly almost-
periodic functions, or Bohr almost-periodic functions, of
Stepanov almost-periodic functions, of Weyl a1most-
periodic functions and of Besicovitch almost-periodic
functions. These classes of almost-periodic functions are
invariant under addition. Together with j(x), each class
also contained the functions f(x), I j(x) I and j(x )e iAx ,
where A is a real number. The metrics DSj[f(x), cp(x)]
are topologically equivalent for all values of I; therefore
it may be assumed that 1= 1. Let SI( -a.p. = SP -a.p., S 1_
a.p.=S-a.p., and BI-a.p.=B-a.p .. Then
U-a.p. C SP-a.p.cWp-a.p. C BP-a.p., p~l.
If PI <P2 and p 1 ;;;'1, then
Sp, -a.p. C SPI -a.p., Wi" -a.p. C WPI -a.p.,
BP' - a.p. C BPI - a.p ..
For every j(x)EB-a.p., the mean value
1 T
M{f(x)} = lim -T jf(x)dx
0
T-HXl
exists. The function a(A,j)=M{f(x)e- iAx }, where A is
a real number, differs from zero only on a countable
set of values of A; any enumeration of this set is called
the sequence {Ad, k = 1, 2, ... , of Fourier exponents
of j(x).
The numbers AA. =a(Abj) are called the Fourier
coefficients of j(x). With a function j(x) in any of the
classes defined above one can associate its Fourier
series
f(x) '" b4;\,/A..x.
k Fourier series if the Fourier exponents are linearly
For j(x)EB 2 -a.p. one has the Parseval equality
M{lf(x) 12} = ~IAA.. 12.
k The importance of criteria for uniform convergence
The Riesz - Fischer theorem can be generalized to the
class BP-a.p.: Let {Ad, k = 1,2, ... , be arbitrary real
numbers, and let {Ad, k = 1, 2, ... , be complex
numbers for which ~:=I I Ak I <00. Then there is an
j(x)EB 2 -a.p. which has the trigonometric series
~kAkeiA.X as its Fourier series.
There is also a uniqueness theorem: If two functions
j(x)EHG(T) and cp(x)EHG(T) have the same Fourier
series, then
DGlf(x), cjl(x)] = o.
In particular, for uniformly almost-periodic functions
the uniqueness theorem states that j(x)=cp(x) (for
Stepanov almost-periodic periodic functions: almost-
everywhere). A uniqueness theorem in the same sense
as for Fourier- Lebesgue series of 2'17-periodic func-
tions does not hold for Weyl or Besicovitch almost-
periodic functions.
The classes of uniformly almost-periodic and of
ALMOST-PERIODIC FUNCTION
Stepanov almost-periodic functions are, respectively,
non-trivial extensions of the class of continuous 2'17-
periodic functions on R and the class of 2'17-periodic
integrable functions on the interval [0, 2'17]. For these
classes of almost-periodic functions the uniqueness
theorem remains valid.
A consequence of the definition of the classes of
almost-periodic functions through the concept of clo-
sure is the approximation theorem: For every almost-
periodic function j(x) from U (or SP or WJ') and every
>0 there is a finite trigonometric polynomial P(x) in
T, satisfying the inequality
Dulf(x),p(x)] <
(Dsplf(x), P(x)] < , DW"lf(x), P(x)] < f).
The approximation theorem may serve as a starting
point of the definition of various classes of almost-
periodic functions. The approximating polynomials
P(x) may contain 'extraneous' exponents, i.e.
exponents different from the Fourier exponents of j(x).
However, important for some applications of the
approximation theorem is the fact that the exponents
different from the Fourier exponents of j(x) can be
avoided in P(x).
In connection with the representability of almost-
periodic functions by generalized Fourier series, the
problem of convergence criteria for these series arises
and various summation methods for generalized
Fourier series (the Bochner- Fejer method, etc.)
become meaningful. Thus, the following criteria have
been obtained: absolute convergence of a generalized
independent; uniform convergence of a Fourier series
when I Ak I~oo as k~oo or when Ak~O as k~oo.
in the theory of almost-periodic functions is
emphasized by the following theorem: If a tri-
gonometric series ~kakeiAkX is uniformly convergent on
the entire real line, then it is the Fourier series of its
sum S(X)E U-a.p .. Corollary: There exists uniformly
almost-periodic functions with an arbitrary countable
set of Fourier exponents. If particular, the Fourier
exponents of a uniformly almost-periodic function may
have finite limit points or may even be everywhere
dense.
Other definitions of almost-periodic functions of the
above classes rely on the concept of an almost-period
and generalizations thereof.
Besides the concept of closure or that of an almost-
period, the concept of a translation can also be used for
the definition of almost-periodic functions. Thus, a
function j(x) is uniformly almost-periodic if and only if
every infinite sequence of functions
155
ALMOST-PERIODIC FUNCTION
f(x + h \ ),f(x + h 2 ), .. , where the translation
numbers h \, h2' ... , are arbitrary real numbers, con-
tains a uniformly convergent subsequence. This defini-
tion serves as a starting point in considering almost-
periodic functions on groups.
The main results in the theory of almost-periodic
functions remain valid if one considers the concept of a
generalized translation. Other generalizations are possi-
ble and useful: almost-periodic functions with values in
an n-dimensional space or in a Banach or metric space,
and analytic or harmonic almost-periodic functions.
References
[I] BOHR, H.: Almost-periodicfunctions, Chelsea, reprint, 1947
(translated from the German).
[2] BESICOVITCH, A.S.: Almost periodic functions, Cambridge Univ.
Press, 1932.
[3] LEVITAN, B.M.: Almost-periodic functions, Moscow, 1953 (in
Russian).
[4] KUPTSOV, N.P.: 'Direct and converse theorems of approxima-
tion theory and semigroups of operators', Russian Math. Sur-
veys 32, no. 4 (1968), liS-In (Uspekhi Mat. Nauk 23, no. 4
(1968), 117-178)
[5] RUDIN, W.: Fourier analySiS on groups, Benjamin, 1962.
[6] LEVITAN, B.M.: Generalized shift operators and some of their
applications, Moscow, 1962 (in Russian).
[7] KRASNOSEL'SKJi, A.M., BURD, V.SR. and KOLESOV, Yu.S.:
Non-linear almost-periodic oscillations, New York, 1973
(translated from the Russian).
E.A. Bredikhina
Editorial comments. The theory of almost-periodic func-
tions was initiated by H. Bohr, who developed the notion of
a uniformly almost-periodic function in his study of Dirichlet
series. This original definition was the one using (-almost-
periods and the approximation theorem was one of the main
achievements of the theory. For a modern approach see
[A4], Sect. 5 and [A10], Chapt. 1,6. As to the equivalence
of the approach starting from a certain structural property
that is a generalization of pure periodicity, and on the other
hand, the approach starting from approximation by tri-
gonometric polynomials (also in the definition of the various
classes of generalized almost-periodic functions), see [2],
Chapt. 2. An interesting generalization not mentioned above
are the Levitan almost-periodic functions (see [A5], where
they are called N-almost periodiC functions). A new
approach to the theory of uniformly almost-periodic func-
tions is given in [A 1); this has lead to the study of so-called
almost-automorphic functions [A9), a class of functions
closely related to the Levitan almost-periodic functions men-
tioned above [A7).
Outside the realm of harmonic analysis an important
application of (uniformly) almost-periodic functions lies In
the theory of differential equations; see e.g. [7], [A8], [AS],
[A2], Chapt. 4, and [A10], Chapt. 4, S. In this context the
theory of dynamical systems (cf. Dynamical system) IS
relevant and, in particular, the study of various types of
almost-penodic and (or) recurrent motions, see [A2), [A3)
and [A6) (beware of conflicting terminology In the literature).
As a replacement for [6) , [A 11) may be used, which IS of
the same flavour as [6].
156
References
[A1] BOCHNER, S.: 'A new approach to almost periodicity', Proc.
Nat. Acad. Sci. USA 48 (1962), 2039-2043.
[A2] BRONSTETN, I.U.: Extensions of minimal transformation
groups, Sijthoff & Noordhoff, 1979 (translated from the Rus-
sian).
[A3] GoTTSCHALK, W.H. and HEDLUND, G.A.: Topological dynam-
ics, Amer. Math. Soc., 1955.
[A4] KATZNELSON, Y.: An introduction to harmonic analysis,
Dover, reprint, 1968.
[A5] LEVITAN, B.M. and ZHIKOV, V.V.: Almost periodiC functions
and differential equations, Cambridge Univ. Press, 1968
(translated from the Russian).
[A6] NEMYTSKJT, V.V. and STEPANOV, V.V.: Quantitative theory of
differential equations, Princeton Univ. Press, 1960 (translated
from the Russian).
[A7] REICH, A.: 'Prakompakte Gruppen und Fastperiodicitat',
Math. Z. 116 (1970), 216-234.
[A8] SELL, G.R.: Topological dynamics and ordinary differential
equations, v. Nostrand Reinhold, 1971.
[A9] VEECH, W.A.: 'Almost automorphic functions on groups',
Amer. J. Math. 87 (1965),719-751.
[A 10] CORDUNEANU, c.: Almost periodic functions, Interscience,
1961.
[A 11] LEVITAN, B.M.: 'The application of generalized displacement
operators to linear differential equations of the second order',
Amer. Math. Soc. Transl. Series 1 (1950),408-541.
(Uspekhi Mat. Nauk 4, no. 1 (29) (1949), 3-112)
AMS 1980 Subject Classification: 42A 75, 43A60
ALMOST-PERIODIC FUNCTION ON A GROUP - A
generalization of almost-periodic functions defined on
R. Let G be an (abstract) group. A bounded complex-
valued function f (x), X E G, is called a right almost-
periodic function if the family f(xa), where a runs
through the entire group G, is (relatively) compact in
the topology of uniform convergence on G, i.e. if every
sequence of functions f(ax\),f(xa2), ... , contains a
subsequence which is uniformly convergent on G. A left
almost-periodic function on G is defined similarly. It
turns out that every right (left) almost-periodic function
f is also left (right) almost-periodic, and the family
f(axb), where a and b independently run through G, is
(relatively) compact. The latter property is often taken
as a definition of almost-periodic functions on G. The
set of all almost-periodic functions on G is a Banach
space with a norm II f II = SUPXEG I f(x) I
The theory of almost-periodic functions on a group
depends essentially on the mean-value theorem (cf. [5],
[8]). A linear functional Mx {f(x)} defined on the space
of almost-periodic functions is called a mean value if
I) Mx {1}=I, Mx{f(x)}~O for f(x)~O and
Mx{f(x)} >0 for j(x)~O,f=O;
2) Mx {f(xa)} = Mx {f(ax)} =Mx{f(x -I)} =Mx {f(x)} ,
for all a EG.
A unitary matrix function g(x)={glj(X)}i'J=I'
defined on G, is called a unitary representation of G if
g(e)=Ir (e is the identity element of G and Ir is the
identity matrix of order r) and if for all x,y E G,
g(xy) = g(x )g(y). The number r is called the dimension

of the representation g. The matrix entries gij(x) are
almost-periodic functions on G. In the theory of
almost-periodic functions on a group they play the
same role as the functions exp(iA(x)) in the theory of
almost-periodic functions on R.
Two representations g(x) and g' (x) are said to be
equivalent if a constant matrix A exists such that
g'(x)=A -lg(x)A. A representation g is said to be
irreducible if the family of the matrices g(x), x E G, does
not admit a common non-trivial subspace in R'. The
set of all irreducible unitary representations is parti-
tioned into classes of mutually-equivalent representa-
tions. Let one representation be chosen from each
equivalence class and let the set thus obtained be
denoted by S. Then the set of almost-periodic functions
H = {</>x(x)} = {</>x: </>x =gt, gES}
on G turns out to be an orthogonal (though, in general,
uncountable) system with respect to the mean value.
Theorem 1 (the Parseval equality). For an almost-
periodic function f(x) the following equality holds:
M (If(x) 12} = ~ I Mx{f(x~(x)} 12
x Mx{ I </>x(x) 12}
Thus, for only countably many A, Mx{f(x)<I>;\(x)}
differs from zero; the series
is called the Fourier series of f
A representation gES is said to occur in the Fourier
series of an almost-periodic function f if
Mx{f(x)gij(x)}=t=O for some i,j, l";;;;i,j";;;;r.
Theorem 2 (the approximation theorem). The set H is
dense in the space of almost-periodic functions
equipped with the norm
II !II = ~~g I f(x) I,
and every almost-periodic function can be arbitrarily
well approximated by a finite linear combination of
matrix entries of representations occurring in its
Fourier series.
If G is a topological group, then to the definition of
an almost-periodic function should be added the
requirement of its continuity. In this case, the represen-
tations occurring in its Fourier series are also continu-
ous.
If G is an Abelian group, then the continuous unitary
representations are one-dimensional. They are called
the characters of G. The characters of G are denoted by
X and Parseval's equality reads as follows:
Mx{ I f(x) 12} = ~ I an 12, an = Mx{f(x)x,,(x)}.
In the case G = Rn the continuous characters are the
ALMOST-PERIODIC FUNCTION ON A GROUP
functions x(x) = exp(iA'x), where AERn,
A'X =A1Xl + ... +AnXn. Theorems I and 2 imply the
main results in the theory of almost-periodic functions
of a single or of several variables.
The proof of the main statements in the theory of
almost-periodic functions is based on the consideration
of integral equations on a group (cf. [2]). The existence
of sufficiently many linear representations of compact
Lie groups has been proved [3]. In this case, invariant
integration (and consequently, the mean) can be esta-
blished directly. Invariant integration on an abstract
compact group has been constructed [4] depending on
an extension of the Peter - Weyl theory to this case.
The theory of almost-periodic functions on a group
can be deduced (d. [3]) from the Peter- Weyl theory
in the following way. Let f be an almost-periodic func-
tion on a group G and let
p{x,y) = a,sup I f(axb)- f(ayb) I.
bEG
Then the set E={tEG: p(t,e)=O} is a normal sub-
group of G, p is an invariant metric on the quotient
group G / E and f is uniformly continuous on G / E.
The almost-periodicity of f implies that the comple-
tion of G / E in the metric p is a compact group and
Theorems I and 2 follow from the Peter- Weyl theory.
References
[I] LEVITAN, B.M.: Almost-periodicfunctions, Moscow, 1953,
Chapt. 6 (in Russian).
[2] WEYL, H.: 'Integralgleichungen und fastperiodische Funk-
tionen', Math. Ann. en (1927),338-356.
[3] PETER, F. and WEYL, H.: 'Die Vollstiindigkeit der primitiven
Darstellungen einer geschIossener kontinuierlichen Gruppe',
Math. Ann. en (1927),737-755.
[4] NEUMANN, 1. VON: 'Zurn Haarschen Mass in topologischen
Gruppen', Compositio Math. 1 (1934), 106-114.
[5] NEUMANN, 1. VON: 'Almost periodic functions in a group 1',
Trans. Amer. Math. Soc. 36 (1934),445-492.
[6] WElL, A.: c.R. A cad. Sci. Paris Ser. I Math. 200 (1935), 38-40.
[7] YOSIDA, K.: Functional analysis, Springer, 1980, Chapt. 8,
Sect. 4; 5.
[8] MAAK, W.: Fastperiodische Funktionen, Berlin, 1950.
V. V. Zhikov
B.M. Levitan
Editorial comments. Instead of the term 'mean value'
one often uses the term 'invariant mean' (cf. [Al], Sect. 18).
For an Abelian group G the uniformly almost-periodic
functions are precisely those that can be continuously
extended to the Bohr compactification of G.
A unified account of the theory of almost-periodic func-
tions on groups can also be found in [A2] and [A3] , Sect.
41. The basic observation is that the Banach algebra of
(continuous) almost-periodic functions on a (topological)
group G is isomorphic to the Banach algebra of all continu-
ous functions on the so-called Bohr compactification Gc of
G. In this way the theory is reduced to the theory of con-
tinuous functions on a compact group (e.g., the mean-value
theorem corresponds to the normalized Haar measure on
Gc the approximation theorem is nothing else than the
157
ALMOST-PERIODIC FUNCTION ON A GROUP
well-known Peter-Weyl theorem for compact groups, etc.).
The Bohr compactification of G can be characterized as the
reflection of G in the subcategory of all compact groups. By
considering reflections in other subcategories of the
category of all topological groups (or even of all semi-
topological semi-groups) one can define other classes of
almost-periodic functions on groups (or semi-groups), see
[A4]. Weakly almost-periodic functions are of particular
interest in functional-analytic applications (semi-groups of
operators). See also [7] and [AS].
References
[A1] HEWITT, E. and Ross, K.A.: Abstract harmonic analysis, 1,
Springer, 1979.
[A2] WElL, A.: I'lntegration dans les groupes topologiques et ses
applications, Hermann, 1940.
[A3] LooMIs, L.H.: An introduction to abstract harmonic analysis,
v. Nostrand, 1953.
[A4] BERGLUND, J.F., JUNGREN, H.D. and MILNES, P.: Compact
right topological semigroups and generalizations of almost
periodicity, Lect. Notes in Math., Vol. 663, Springer, 1978.
[A5] BURCKEL, R.B.: Weakly almost periodic functions on semi-
groups, Gordon and Breach, 1970.
[A6] CORDUNEANU, c.: Almost periodic functions, Interscience,
1961, Chapt. 7.
[A7] GUCKSBERG, I. and LEEUW, K. DE: 'Almost periodic functions
on semigroups', Acta Math. 105 (1961),99-140.
[A8] AMERIo, L. and PROUSE, G.: Almost-periodic functions and
functional equations, v. Nostrand, 1971.
[A9] DIXMIER, J.: C' algebras, North-Holland, 1977 (translated
from the French).
AMS 1980 Subject Classification: 43A60, 42A75
ALMOST-PRIME NUMBER - A natural number n
of the form
n=PI"'Pk,
where the Pi are prime numbers and k ~ I is a constant.
Prime numbers are the special case of almost-prime
numbers for k = 1. There are theorems for almost-prime
numbers that generalize theorems on the distribution of
prime numbers in the set of natural numbers. Several
additive problems that have not yet been solved for
prime numbers have been solved for almost-prime
numbers.
B.M. Bredikhin
Editorial comments. See also Sieve method.
An example of a result on the distribution of almost-prime
numbers generalizing the corresponding one on prime
numbers is the following. Let 'TI'k(X) be the number of
square-free almost-prime numbers .;;;x. Then (cf. [A 1], Sect.
22.18) x(log log xt- 1
'TI'k(X) - (k-2)! logx
References
[A1] HARDY, G.H. and WRIGHT, E.M.: An introduction to the
theory of numbers, Clarendon Press, 1965.
AMS 1980 Subject Classification: 10A25
ALMOST-REDUOBLE LINEAR SYSTEM of ordi-
nary differential equations - A system
x = A(t)x, xERn, (*)
158
A(): R -+ Hom(Rn, Rn),
having the following property: There exist a system
j=By, YERn , with constant coefficients and, for every
>0, a Lyapunov transformation L((t) such that by the
change of variables x =L((t)y, the system (*) is
transformed into the system
j = (B+C.(ty,
where
~~f II C.(t) " < (.
Every reducible linear system is almost reducible.
References
[1) lzoBov, N.A.: 'Linear systems of ordinary differential equa-
tions', J. Soviet Math. S, no. 1 (1976),46-96. (Itogi Nauk. i
Tekhn. Mat. Anal. 12 (1974),71-146) VM M'll' h h'k
llOnSClOV
AMS 1980 Subject Classification: 34A34 , 34A30
ALMOST-SYMPLECTIC STRUCfURE - A non-
degenerate differential 2-form on a manifold. An
almost-symplectic structure 1'2 can exist only on an
even-dimensional manifold M (dimM=2m) and
defines an Sp(m, R)-structure BSp(m,R), namely the prin-
cipal fibre bundle of frames on M with structure group
Sp(m, R), consisting of all frames r = {ei' Ji:
i = 1, ... ,m} for which
Q(e;, ej) = Q(j;,jj) = 0, Q(e;,jj) = 81)'
A necessary and sufficient condition for the existence
of an almost-symplectic structure (or of an almost-
complex structure, as well) on a manifold M is the pos-
sibility of reducing the structure group of the tangent
bundle to the unitary group U(m). For this, in particu-
lar, it is necessary that all odd-dimensional
Stiefel- Whitney classes of M vanish (cf. [I D.
An almost-complex structure J and a Riemannian
metric g on a manifold M define an almost-symplectic
structure Q by the formula
Q(X, Y) = g(JX, Y)-g(X,JY),
where X and Yare vectors. Any almost-symplectic
structure can be obtained in this manner. An almost-
symplectic structure is said to be integrable or, in other
words, a symplectic structure, if it can be brought to the
form Q = ~dxi I\dyi in some local coordinates Xi ,yi,
i = 1, ... , m, in a neighbourhood of any point. Accord-
ing to Darboux's theorem, for this it is necessary and
sufficient that Q be closed. An example of an integrable
almost-symplectic structure is the canonical symplectic
structure Q = ~dpi I\dqi on the cotangent bundle T* M
of an arbitrary manifold M (here the ql are local coor-
dinates on M and the pi are the associated coordinates
in the fibres). An example of a non-integrable almost-
symplectic structure is a left-invariant 2-form on a
semi-simple Lie group G, obtained by extending an

arbitrary non-degenerate exterior 2-form on the Lie
algebra TeG of G by left translation to G. As a
Riemannian metric, an almost-symplectic structure also
defines an isomorphism of the tangent and cotangent
spaces (and by the same method, of the spaces of con-
travariant and covariant tensors); it further defines a
canonical 2m-form 1/ = om / m!, called its volume form,
and several operators in the space /\(M) of differential
forms: the operator n of exterior multiplication by 0;
the operator in of interior multiplication by 0; the
Hodge star operator *: /\p(M)-,>/\2m -P(M), w-'>iw1/,
where the operator i w of interior multiplication is
defined as the contraction of the given form with the
p-vector corresponding to the p-form w; the operator of
codifferentiation 0 = *d*. In contrast with the Riemann-
ian case, the operator !:l=do+od turns out to be skew-
symmetric with respect to the global scalar product
<a, [3> = fMa/\ *[3 in the space of p-forms on a com-
pact manifold M. For an arbitrary p-form one has the
Hodge - Lepage decomposition
W=Wo +nWl +bW2 + .. " where the Wi EN -2i(M)
are uniquely determined effective forms (i.e. they are
annihilated by in) [3].
An almost-symplectic structure is said to be confor-
mally flat if there is a function A>O such that
d(AO) = O. This is equivalent to the representability of 0
in the form
Q = v 1 ~ dx' /\dyl
, - 1
For m =2, a necessary and sufficient condition in order .
that the almost-symplectic structure 0 be conformally AMS 1980 Subject Classification: 03C07
flat is the closedness of the I-form 00 = indO, and for ~
m>2 the equality dO=(l/m-l)oO/\O should hold
(cf. [1]). group An of the symmetric group Sn consisting of all
The tensor T of type (1,2) corresponding to the 3-
form dO and defined by the equality
O(TxY, Z)=dO(X, Y, Z), where X, Yand Z are vectors,
is called the torsion tensor of the almost-symplectic struc-
ture O. The (degenerate) metric g(X, Y)=trTxT y can
be associated with it. An almost-symplectic structure
determines the class of linear connections V for which
o is parallel and which have T as their torsion tensor.
Two such connections differ by a tensor field of the
form Oij Sjkl, where Sjkl is an arbitrary symmetric tensor
field. The connections under consideration correspond
in a one-to-one manner to the sections of the first
extension B l-'>B for the Sp(m, R)-structure
B =B Sp(m. R), which is the principal bundle of frames on
B with structure group s3(R2m) (the vector group of
homogeneous polynomials in 2m variables of degree 3).
The Sp(m, R)-structure is a G-structure of infinite type.
Therefore, the group of automorphisms of an almost-
symplectic structure can be infinite-dimensional. In
particular, the group of automorphisms of a symplectic
ALTERNATING KNOTS AND LINKS
structure is always infinite-dimensional and is a k-
transitive group for any k >0.
References
[I) LIBERMAN, P.: 'Sur les structures presque complexe et autres
structures infinitesimales regu!ieres', Bull. Soc. Math. France 83
(1955), 195-224.
(2) Itogi Nauk i Tekhn. Algebra Topol. Geom. 11 (1974), 153-207.
(3) LYCHAGIN, V.V.: 'Contact geometry and second-order non-
linear differential equations', Russian Math. Surveys 34, no. I
(1979), 149-180. (Uspekhi Mat. Nauk 34, no. I (1979), 137-
165)
(4) KOBAYASHI, S.: Transformation groups in differential geometry,
Springer, 1972.
(5) HURT, N.E.: Geometric quantization in action, Reidel, 1983.
(6) ARNOL'D, V.1. and GIVENTAL, A.B.: Itogi Nauk. i Tekhn.
Sovrem. Probl. Mat. 4, 5-139.
D. V. Alekseevskir
Editorial comments.
References
[A 1] LIBERMANN, P. and MARLE, C.-M.: Symplectic geometry and
analytical mechaniCS, Reidel, 1987.
AMS 1980 Subject Classification: 53C15
ALPHABET - A list (finite set) of letters (cf. Letter).
An alphabet usually contains the letters which are
required to develop a certain system of symbols or, as
is sometimes said, a language. For two alphabets A and
B, their union A U B, their intersection A n
B, their
difference A \ B, as well as the inclusion relation A CB
are defined in the obvious manner. For convenience,
the concept of an empty alphabet, i.e. an alphabet
without any letters, may also be considered.
NM. Nagornyr
ALTERNATING GROUP of degree n - A The sub-
even permutations. An is a normal subgroup in Sn of
index 2 and order n! / 2. The permutations of An, con-
sidered as permutations of the indices of variables
x I, . . . ,Xn , leave the alternating polynomial
II(xi-Xj) invariant, hence the term 'alternating
group'. The group Am may also be defined for infinite
cardinal numbers m, as the subgroup of Sn consisting
of all even permutations. If n > 3, the group An is
(n - 2)-fold transitive. For any n, finite or infinite,
except n =4, this group is simple; this fact plays an
important role in the theory of solvability of algebraic
equations by radicals.
References
[I) HALL, M.: The theory of groups, Macmillan, 1959.
NN Vi/'yams
Editorial comments. Note that A5 is the non-Abelian
simple group of smallest possible order.
AMS 1980 Subject Classification: 20835
ALTERNATING KNOTS AND LINKS - Knots and
159
ALTERNATING KNOTS AND LINKS
links with an alternating diagram (cf. Knot and link
diagrams), i.e. a projection in general position onto a
plane such that, when successively traversing all com-
ponents, the overpasses and underpasses succeed each
other in alternation. Any diagram can be converted to
an alternating diagram by changing the upper and
lower branches through its double points.
Let F be a Seifert surface. As distinct from the ordi-
nary case, the inequality d.,:;;;2h+p.-l, where d is the
degree of the Alexander polynomial (d. Alexander
invariants), h is the genus of the Seifert surface and p. is
the number of components of the link k, becomes an
equality in the case of alternating knots and links.
Accordingly, the genus of an alternating link can be
computed from anyone of its alternating diagrams, the
Seifert surface being the surface of smallest genus. This
also shows that if the diagram is normalized, i.e. if the
projection plane does not contain a simple closed con-
tour that intersects the diagram at one double point,
the link is trivial (cf. Knot theory) if and only if the
diagram contains no double points. If such a contour
exists, it is possible, by rotating an internal part of the
diagram through 180 degrees, to reduce the number of
double points while preserving the alternating nature of
the diagram. This yields an algorithm for solving the
problem of triviality of alternating knots and links.
Moreover, if the diagram is connected, the links do not
become separated, since d;;;;.l, and the reduced Alex-
ander polynomial of seperated links is zero. The Alex-
ander matrix is computed as the incidence matrix of
some graph, which implies [1], [2] that Ll(t) is an alter-
nating polynomial, i.e. its coefficients are non-zero and
their signs alternate. If Ll(O) = I, alternating knots and
links are known as Neuwirth knots and links (d.
Neuwirth knot). The number of double points of a nor-
malized diagram of an alternating knot or link is not
larger than its determinant. The groups of alternating
knots and links (d. Knot and link groups) are
represented as free products with amalgamation of two
free groups of rank q by subgroups of rank 2q - 1. This
representation is obtained with the aid of the van Kam-
pen theorem, if the space of the link k is subdivided by
the boundaries of regular neighbourhoods of the Seifert
surface with respect to k, constructed over the alternat-
ing diagram. All knots of the standard table (cf. Knot
table) with non-alternating diagrams are non-
alternating knots. Most parallel knots, windings, etc.,
do not alternate.
References
[I] MURASUGI, K.: 'On the Alexander polynomial of the alternat-
ing knot', Osaka J. Math. 10 (1958), 181-189.
[2] CROWELL, R.H.: 'Genus of alternating link types', Ann. of
Math. (2) 69 (1959), 258-275.
[3] MURASUGI, K.: 'On alternating knots', Osaka 1. Math. 12
(1960),277-303.
A. V Chernavskir
160
Editorial comments. A Seifert surface for a knot or link
K n c sn+2 is a connected, bicoliared compact manifold
M n + 1 CS n + 2 such that aM n + 1 =Kn. A subset Xc Yis bicof-
fared (in Y) if there exists an imbedding b: XX [ -1, 1]~ Y
such that b(x, O)=x for ali XEX. The mapping b or its image
is the bicoliar itself.
References
[A 1] ROLFSEN, D.: Knots and links, Publish or Perish, 1976.
AMS 1980 Subject Classification: 57M25, 57Q45
ALTERNATING SERIES - An infinite series whose
terms are alternately positive and negative:
U]-U2+ ... +(-lr-]un + " ' , Uk> O.
If the terms of an alternating series are monotone
decreasing (un + 1 <Un) and tend to zero (limn->oo Un =0),
then the series is convergent (Leibnitz' theorem). The
remainder term of an alternating series,
'n = (-ltun +]+"',
has the same sign as its first term and is less then the
latter in absolute value. The simplest examples of alter-
nating series are
1-1-+1--1-+ ... +(_l)n-]1-+ ...
2 3 4 n'
1-1-+1--1-+ ... +(-lr-]_I_+ ....
2 5 7 2n-l
The sum of the first of these series is log 2; that of the
second is 'IT / 4.
v,l. Bityutskov
AMS 1980 Subject Classification: 40A05
ALTERNATION, skew symmetry, anti-symmetry,
alternance - One of the operations of tensor algebra,
yielding a tensor that is skew-symmetric (over a group
of indices) from a given tensor. Alternation is always
effected over a few superscripts or over a few sub-
scripts. A tensor A with components {aj',... J" I ~i p,
j p..,:;;;n} is the result of alternation of a tensor T with
components {tj', ...... J,' I":;;;i p , jp..,:;;;n}, for example, over
superscripts, over a group of indices 1= (i 1, . . . ,im ) if
(*)
The summation is conducted over all m! rearrange-
ments (permutations) a=(a], ... ,am) of I, the number
a(I, a) being + 1 or -1, depending on whether the
respective rearrangement is even or odd. Alternation
over a group of subscripts is defined in a similar
manner.
Alternation over a group of indices is denoted by
enclosing the indices between square brackets. Secon-
dary indices inside the square brackets are separated by
vertical strokes. For instance:

Successive alternation over groups of indices I I and h,
I I Ch, coincides with alternation over the group of
indices 1 2 :
If n is the dimension of the vector space on which
the tensor is defined, alternation by a group of indices
the number of which is larger than n will always pro-
duce the zero tensor. Alternation over a given group of
indices of a tensor which is symmetric with respect to
this group (d. Symmetrization (of tensors also yields
the zero tensor. A tensor that remains unchanged under
alternation over a given group of indices I is called
skew-symmetric or alternating over I. Interchanging any
pair of such indices changes the sign of the component
of the tensor.
The operation of tensor alternation, together with the
operation of symmetrization, is employed to decompose
a tensor into simpler tensors.
The product of two tensors with subsequent alterna-
tion over all indices is called an alternated product
(exterior product).
Alternation is also employed to produce sign-
alternating sums of the form (*) with multi-indexed
terms. For instance, a determinant with elements which
commute under multiplication can be computed by the
formulas
al
a7
References
[1] SHIROKOY, P.A.: Tensor calculus. Tensor algebra, Kazan', 1961
(in Russian).
[2] BEKLEMISHEY, D.V.: A course of analytical geometry and linear
algebra, Moscow, 1971 (in Russian).
[3] SCHOlITEN, 1.A.: Tensor analysis for physicists, Cambridge
Uniy. Press, 1951.
[4] EFIMOY, N.V. and ROZENDORN, E.R.: Linear algebra and
multi-dimensional geometry, Moscow, 1970 (in Russian).
L.P. Kuptsov
AMS 1980 Subject Classification: 15A69, 15A72,
53A45
ALTERNATION, POINTS OF - A sequence of points
{x 1 }3+ 1 C Q, a";;xO<xl<'" <xn +l.,;;b,
at which the difference (Xi = f(x;)- Pn(x;) assumes non-
zero values of alternating signs. Here f (x) is a con-
tinuous function on the closed set Q qa, bland Pn(x)
is an algebraic polynomial of degree not exceeding n.
In a similar manner the concept of points of alternation
is introduced for polynomials in a Chebyshev system of
functions {sk(x)}8 (which satisfy the Haar condition).
If, in this situation, all absolute values of (Xi are equal
to
ALTERNATIVE RINGS AND ALGEBRAS
max If(x)-Pn(x) I,
XEQ
then the points {xd8 + I are called Chebyshev points of
alternation. Points of alternation play an important role
in the theory of approximation of functions. E.g., the
de la Vallee-Poussin theorem (the alternation theorem)
and the Chebyshev criterion (d. Chebyshev alternation)
are formulated in terms of points of alternation. Points
of alternation are also employed in constructing poly-
nomials of best approximation.
Yu.N. Subbotin
Editorial comments. A sequence of Chebyshev points of
alternation is also called an alternating set [A 1], Chapt. 1.
References
[A 1] T.1.: An introduction to the approximation of func-
RlvLIN,
tions, Dover, reprint, 1981.
[A2] MOLLER, M.W.: Approximationstheorie, Akad. Ver-
lagsgesellschaft, Wiesbaden, 1978.
[A3] MEiNARDus, G.W.: Approximation von Funktionen und ihre
numerische Behandlung, Springer, 1964.
[A4] CHENEY, E.W.: Introduction to approximation theory, Chel-
sea, reprint, 1982.
AMS 1980 Subject Classification: 41 A 10
ALTERNATIVE (in game theory) - Each of the posi-
tions to which, according to the rules of the game, it is
possible to pass from a given position in one move. The
move itself, and the index of the move, are also known
as alternatives.
I.N. Vrublevskaya
AMS 1980 Subject Classification: 90DXX
ALTERNATIVE RINGS AND ALGEBRAS - An alter-
native ring is a ring in which every two elements gen-
erate an associative subring; an alternative algebra is a
(linear) algebra that is an alternative ring. By a theorem
of E. Artin the class of all alternative rings is defined
by the system of identities:
(xy)y = x(yy) (right altemativeness);
(xx)y = x(xy) (left altemativeness).
Thus, the alternative rings form a variety. The termi-
nology 'alternative ring' is justified by the fact that in
any such ring the associator (the defect of associativity)
(x,y,z) = (xy)z-x(yz)
IS a skew-symmetric (alternative) function of its argu-
ments.
The first example of alternative rings were the Cayley
numbers, which form an alternative skew-field, i.e. an
alternative ring with a unit element, in which the equa-
tions ax = band ya = b are uniquely solvable for all b
and all a*O. Alternative skew-fields playa substantial
part in the theory of projective planes, since a projec-
tive plane is a Moufang plane (i.e. a translation plane
with respect to some straight line) if and only if any
coordinatization of its ternary ring is an alternative
skew-field. If in a ring R with a unit element each
161
ALTERNATIVE RINGS AND ALGEBRAS
non-zero element is invertible, and if the identity
a-l(ab) =b
(or the identity (ba)a -) =b) is valid for any a, b ER,
then R is an alternative skew-field. Any alternative
skew-field is either associative or has the structure of a
Cayley - Dickson algebra over its centre.
Each simple alternative ring is also either associative
or a Cayley- Dickson algebra over its centre (in this
case the algebra need not necessarily be a skew-field).
Both the associative and the primitive alternative rings
are exhausted by the Cayley - Dickson algebras. All
primary alternative rings R (if 3R:~0) are either associ-
ative or are Cayley- Dickson rings.
Many properties of alternative rings differ substan-
tially from those of an associative ring in a similar
situation. Thus, if R is an alternative ring and A and B
are right ideals in it, their product AB need not neces-
sarily be a right ideal, even if A is a two-sided ideal in
R; however, the product of two two-sided ideals of an
alternative ring is a two-sided ideal in it. The differ-
ence between associative rings and alternative rings is
also strongly manifested by the fact that alternative
rings contain various kinds of nilpotency, since the pro-
duct of elements may be zero or non-zero, depending
on the placement of the parentheses. It is customary to
use the following kind of nilpotency in an alternative
ring: solvability (a ring R is called a solvable ring of
index m if there exists a number m such that Rm = 0,
where R; + I =R; R;, R) = R), right nilpotency (there
exists a number n such that R(n) =0, where
R(i+l)=R(i)R(l), R(I)=R), and nilpotency (there exists
a number k such that Rk =0, i.e. the product of any k
elements of R is zero, no matter how the parentheses
are placed). There are solvable alternative rings of
index 3, but they are not nilpotent. Right nilpotency in
an alternative ring is equivalent to nilpotency (alterna-
tive rings that are right nilpotent of index n, are also
nilpotent of index ~(n + 1)2). Locally, i.e. in finitely
generated rings, all kind of nil potency are equivalent.
The theory that establishes sufficient criteria of local
nilpotency of an alternative ring is completely parallel
to the corresponding theory for associative rings. This
is a consequence of the following fact: Let R be an
alternative ring in which one can to select a system of
generators S such that any two elements of S generate
a nil ring, suppose also that all assocIative
homomorphic images of R are locally nilpotent. then R
is locally nilpotent. Therefore. if R is an alternative
ring with an identity xn =0. then R is locally nilpotent;
if R is an alternative algebra with an identity relation
that is not a consequence of the associativity. and if
each element is the sum of a finite numher of nil ele-
ments. then R is locally nilpotent. As regards global. as
162
distinct from local, nilpotency the situation in alterna-
tive rings differs from that in associative rings. Thus, an
alternative ring with the identity x 3 =0 is not neces-
sarily nilpotent (even if its additive group is torsion-
free). However, an alternative ring with an identity
xn =0 and without elements of order k, O<k~n, in the
additive group is solvable of index n (n + I) / 2.
If R is an algebraic alternative algebra with an iden-
tity relation that is not a consequence of the associa-
tivity (or if the degrees of algebraicity of the elements
of R are uniformly bounded), then R is locally finite-
dimensional.
In a alternative ring there is an analogue of the
Jacobson radical: In each alternative ring there exists a
maximal quasi-regular ideal J (R), which is equal to the
intersection of all modular maximal right ideals. The
quotient ring R / J (R) is J -semi-simple, l.e.
J(R / J(R=O; if I is an ideal of R, then
J(I)=J(R)nI, and any J-semi-simple ring can be
approximated by primitive alternative rings (i.e. by
primitive associative rings and by Cayley - Dickson
algebras). There also exist analogues of all other associ-
ative radicals (lower nil-radical, locally nilpotent radi-
cal, etc.), which display the same basic properties as in
associative rings.
In an Artinian alternative ring R (i.e. an alternative
ring that satisfies the descending chain condition
(minimum condition) for right ideals) the radical R (J)
is nilpotent, and so is any nil-sub groupoid of the multi-
plicative groupoid of R. A ring R is an Artinian alter-
native ring without nilpotent ideals if and only if R
splits into the direct sum of finitely many complete
matrix algebras (over certain associative skew-fields)
and Cayley- Dickson algebras; this decomposition is
unique for each R, apart from a rearrangement of the
terms. If R is an alternative ring, lone of its ideals,
and R satisfies the minimum condition for the two-
sided ideals contained in I, then I is nilpotent if and
only if there are no simple ideals in the ring Ref> for any
arbitrary homomorphism ef> from R into Ief>.
In an alternative ring R one distinguishes between
the associative centre
N(R) = {nER: (n,a,b)=Oforalla,bER},
the commutative centre
C(R) = {cER: [c,a]=ca-ac=OforallaER},
and the centre
Z(R) = N(R) n C(R).
If in the additive group of R there are no elements of
order three. then C(R)CN(R). However. commutative
non-associative alternative algebras exist over a field of
characteristic three. In a primary alternative ring R one
always has C (R) ~ N (R). In any alternative ring R one
 AL TERNA TIVE RINGS AND ALGEBRAS

always has [N (R), R] C;N (R). Suppose that in an alter-
native ring R there are no non-trivial ideals. Then 1)
either 3R =0 or N(R)=!=O; 2) either 3R C;N(R) or
Z (R)=!=O; 3) if A is a right ideal in R, then
N(A)=A nN(R) and Z(A)=A nZ(R). However,
[7) SKORNYAKOY, L.A.: 'Alternative rings', Rend di Mat. e Appl.
24 (1965), 360-372.
[8) SLATER, M.: 'Ideals in semiprime alternative rings', J. of Alge-
bra 8 (1968), 60-76.
[9) DoROFEEV, G.V.: 'Alternative rings with three generators',
Sibirsk. Mat. Zh. 4, no. 5 (1963), 1029-1048 (in Russian).

over any field F it is possible to construct an alterna- KA. Zhevlakov

tive ring K with no non-trivial ideals such that Editorial comments. An element z in a ring R is right
N(K)=C(K)=O, but K2, which is an ideal of K, is an quasi-regular if {a - za: a E R} = R. An element is quasi-
associative-commutative ring, i.e. regular if it is both left and right quasi-regular. An ideal in R

N(K2) = C(K2) = K2 * O.
is quasi-regular if all its elements are quasi-regular [A2). The
phrase 'approximated by ... ' means 'representable as a
Of the identities that are valid in an associative ring, subdirect sum of ... '. For example 'a J-semi-simple alter-
the following are the best known: native algebra R can be approximated by primitive alterna-
tive rings' means that R is isomorphic to a subdirect sum of
[(xy)z]y = x[(yz)y], primitive alternative rings.
[(xy)x]z = x [y(xz)], References
(xy)(zx) = [x(yz)]x [A 1] JACOBSON, N.: Structure and representation of Jordan alge-
bras, Amer. Math. Soc., 1968.
(the Moufang identities); [A2] JACOBSON, N.: Structure of rings, Amer. Math. Soc., 1964.

(xy, z, I)-y(x, z, I)-(y, z, I)x
= ([x,y), z, 1)+(X,y, [z, Ij),
([x,y]", z, I) = [x,y]([x,yf, z, I) =
= [A3] SCHAFER, R.D.: An introduction to non-associative algebras,
Acad. Press, 1966.
AMS 1980 Subject Classification: 16-XX, 17D05

= ([x,yf, z, 1)[X,y] = o. ALTERNION - A hypercomplex number. Alternions

An alternative ring with three generators also satisfies
the identity
([x,y][z, I]+[z, I][X,y], u, v) = o.
In an alternative ring with more than three generators
the identity (*) is usually not satisfied; moreover, in
these alternative rings, as a rule, ([x, y f, z, t)=!=O. In
any alternative ring without locally nilpotent ideals, the
identity (*) is satisfied, so that such an alternative ring
may be considered as a generalization of the complex
numbers, double numbers (cf. Double and dual
(*) numbers) and quaternions. The algebra IAn of altern ions
of order n and of index I is an algebra of dimension
2n -lover the field of real numbers, with unit element
1 and a system of generators II, ... ,In -I, in which the
multiplication satisfies the formula
li~ = -~/i' If = -(i,

can be approximated by primary associative rings and where fj = + 1, the value - 1 occurs I times and + 1
by Cayley - Dickson rings. occurs n -1-1 times, respectively. A base of the alge-
Any free alternative ring R has a non-zero ideal bra is formed by the unit element and by elements of
U(R) contained in the associative centre N(R). A free the form
alternative ring with three or more generators not only
contains divisors of zero, but is not primary. Free alter- where h < ... <A. In this base any alternion a can
native rings with four or more generators even contain be written as
trivial ideals and for this reason cannot be approxi- a = a+ ~ai/i+ ~~aiJlij+ ... +
mated by primary rings. i ; j

References +a I '(n-I)l I, ... n-I>

[1) DoROFEEY, G.V.: 'An instance of a solvable, though non-
nilpotent alternative ring', Russian Math. Surveys 15, no. 3
where a, ai, ... ,a I ... (n -I) are real numbers. The
(1960), 147-150. (Uspekhi Mat. Nauk 15, no. 3 (1960), 147- alternion a conjugate to the alternion a is defined by
150) the formula
[2) ZHEVLAKOY, K.A.: 'Alternative Artinian rings', Algebra i
Logika 5, no. 3 (1966), 11-36 (in Russian). a = ~(_1)k(k+I)/2ail "'i'/il ... Ii,.
[3) ZHEVLAKOY, K.A.: 'On radicals and von Neumann ideals', k
Algebra and LogiC 8, no. 4 (1969), 244-252. (Algebra i Logika The following equalities hold
8, no. 4 (1969), 425-439)
[4) SHIRSHOY, A.I.: 'On certain non-associative nil-rings and alge- a+f3 = a+P, ~ = a, af3 = pa.
braic algebras', Mat. Sb. 41 (83), no. 3 (1957), 381-394 (in Rus-
sian). The product aa is always a positive real number; the
[5) KLEINFELD, E.: 'Simple alternative rings', Ann. of Math. (2) 58,
no. 3 (1953), 544-547.
quantity I a I = y:;;;;
is called the modulus of the
[6) KLEINFELD, E.: 'Alternative nil rings', Ann. of Math. (2) 66, no. altern ion a. If the number I P- a I is taken as the dis-
3 (1957), 395-399. tance between two alternions a and p, then the alge-

163
ALTERNATIVE RINGS AND ALGEBRAS
bras 0 An and 'An, 1>0, are isometric to the Euclidean
space R 2"-' and the pseudo-Euclidean spaces 'R2"-',
respectively. The algebra 0A 1 is isomorphic to the field
of real numbers; 0 A 2 is isomorphic to the field of com-
plex numbers; 1A 2 is isomorphic to the algebra of dou-
ble numbers; 0 A 3 is isomorphic to the skew-field of
quaternions; and 1A 3 and 2A 3 are isomorphic to the
so-called algebras of anti-quaternions. The elements of
oAn are the so-called Clifford numbers. The algebra
4A 5 was studied by P. Dirac in the context of the spin
of an electron.
The algebras of alternions are special cases of Clif-
ford algebras (d. Gifford algebra).
References
[1) ROZENFEL'D, B.A.: Non-Euclidean geometry, Moscow, 1955 (in
Russian).
N. N. Vi/'yams
AMS 1980 Subject Classification: 16AXX
AMALGAM - A set M represented as the set-
theoretic union of a family {Ai: i EI} of algebraic sys-
tems (cf. Algebraic system) Ai of a given class ~ with
intersections Vi}' where for all i, j the intersection
Ai nAj = Ui] = Uji
is non-empty and is a subsystem of each of the systems
Ai, A j . If there exists a system B in the class Sf contain-
ing all Ai (i EI) as subsystems, then says one that the
amalgam is imbeddable in the system B.
An amalgam of two groups and, in general, any
amalgam of groups {Ai: i EI}, in which all intersec-
tions Vi} (i=l=j) coincide and are equal to V, is always
imbeddable in a group, e.g. in the free product of the
groups Ai (i EI) with the common subgroup U. There
are, however, amalgams of groups which are not
imbeddable in a group. (For conditions for imbeddabil-
ity of amalgams of groups in a group see [l]; for
imbeddability of amalgams of semi-groups in a semi-
group see [2]). See also Amalgam of groups.
Let ~ be the class of all algebras over a given field F
or the class of commutative, anti-commutative or Lie
algebras over a field F. An amalgam {A;: i EI} of Sl-
algebras Ai with identical intersections Vi} = V (for all
i =1=j) is imbeddable in the S'f-free product of these alge-
bras with the common subalgebra V [3]. It has been
shown [4] that an amalgam {Ti: i EI} of associative
skew-fields Ti with identical intersections Tij = T (i=l=j)
is imbeddable in an associative skew-field.
References
[1) KUROSH, A.G.: The theory of groups. 2. Chelsea. reprint. 1960.
33 (translated from the Russian).
[2) CLIFFORD. A.H. and PRESTON. G.B.: The algefTraic theory of
semigroups. 2. Amer. Math. Soc., 1967.
[3) SHIRSHOY, A.I.: 'On a hypothesis in the theory of Lie alge-
bras'. Sibirsk. Mat. Zh. 3, no. 2 (1962). 297-30 I (in Russian).
164
[4] COHN, P.M.: 'The embedding of firs in skew fields', Proc. Lon-
don Math. Soc. (3) 23 (1971), 193-213.
L.A. Bokut'
D.M. Smirnov
Editorial comments. The reference to [1] in the article
above should have been to H. Neumann's original papers
[A 1], [A2].
The amalgamation problem for a class of algebraic sys-
tems is generally understood to mean the problem of imbed-
ding an amalgam of two systems, as defined above, in a
system of the class. As mentioned above, the amalgamation
problem for groups is solvable, but for other classes it can
be an interesting and important question. E.g., it has been
shown [A3] , [A4] that the amalgamation problem for various
classes of lattices is closely related to the interpolation prob-
lem in logic.
References
[A1] NEUMANN, H.: 'Generalized free products with amalgamated
subgroups', Amer. J. Math. 70 (1948), 590-625.
[A2] NEUMANN, H.: 'Generalized free products with amalgamated
subgroups', Amer. J. Math. 71 (1949), 491-540.
[A3] MAKSIMOYA, L.L.: 'Craig's interpolation problem and amal-
gamable varieties', Soviet Math. Oak/. 18 (1977), 1550-1553.
[A4] PrITs, A.M.: 'Amalgamation and interpolation in the category
of Heyting algebras', J. Pure App/. A/g. 29 (1983), 155-165.
AMS 1980 Subject Classification: 08825
AMALGAM OF GROUPS - A family of groups
G a, a EI, that satisfies the condition that the intersec-
tion Ga n G p is a subgroup in G a and G p for any a, f3
from I. An example of an amalgam of groups is an
arbitrary family of subgroups of a given group. An
imbedding of an amalgam of groups A = {G a: a EI} into
a group G is a one-to-one mapping of the union
U aE I Ga into G whose restriction to each Ga is an iso-
morphism. An amalgam of groups in which all intersec-
tions Ga n Gp are identical (and equal to, say, a sub-
group H) is imbeddable in the group that is the free
product of the groups Ga with the amalgamated sub-
group H. On the other hand, there exists an amalgam
of four Abelian groups that is not imbeddable in a
group. The principal problem concerning amalgams of
groups is, generally speaking, as follows. Let a, T be
possible properties of groups. The question to be
answered is the nature of the conditions under which
an amalgam of groups with the property a is imbedd-
able in a group with the property T. It was found that
all amalgams of two finite groups are imbeddable in a
finite group. An amalgam of three Abelian groups is
imbeddable in an Abelian group. An amalgam of four
Abelian groups imbedded in a group is contained in an
Abelian group. There exists an amalgam of five Abelian
groups which is imbeddable in a group, but not in an
Abelian group. Another problem that has been studied
is the imbeddability of an amalgam of groups if a, T
denote solvability, nilpotency, periodicity, local finite-

ness, etc. (in different combinations).
Yu.I. Merzlyakov
N.S. Romanovskii
Editorial comments. In the definition of amalgam above,
think of the Ga as all being subsets of some large set X.
The amalgamated product of groups G; 'over a common
subgroup U' is constructed as follows. Let G; be a set of
groups indexed by the set I. For each i let U; be a sub-
group of G; and for each i let there be an isomorphism
q,;: U;~U identifying U; and U. Consider the set G of all
words
a, ... at
with each a; from some Gj , and consider the following ele-
mentary equivalences
1) if a;=1 then a, ... a;_,a;a;+, ... at is equivalent to
a, ... a;_, a;+, ... at;
2) if a; and a; +' belong to the same group Gj and
aja;+, =a; in Gj then a, ... aj_, a;aj+, ... at is equivalent
to a, ... aj_, a;aj+2 ... at;
3) if aj=UjEUjCGj and bj=UkEUkCGk and
q,k(Uk)=U=q,j(U;)EU, then a, aj_,aja;+, ... at is
equivalent to a, ... aj_, bja;+, ... at.
Let - be the equivalence ~elation generated by these ele-
mentary equivalences, then G / - is the amalgamated pro-
duct of the G;, more precisely of the (Gj, Uj) (Le. the free
product of the G; with amalgamated subgroup ll); the group
law is induced by concatenation.
Amalgamated products are non-trivial. This follows from
the following canonical form theorem. For each i select a
set R; of left coset representatives of U in G;. Then each
word is equivalent to precisely one of the form UZ, ... Zt
with each Z; in some Rj , UEU and z, and Zj+, belonging to
different G/s for i=1, ... ,t-1. If U={e} one obtains of
course the free product of the Gi . A subgroup of a free pro-
duct is itself a free product (Kurosh' theorem). Subgroups of
a product with an amalgamated subgroup need not be
themselves of this type. The reason is that if U is the amal-
gamated subgroup, then one can take subgroups Hj of the
Gj with different intersections with U so that the Hj will amal-
gamate in various different ways. This leads to generalized
amalgamated products and the notion of amalgam as
defined above. The theory of these is still incomplete.
References
[A1] HALL, M, JR.: The theory of groups, Macmillan, 1959.
[A2] NEUMANN, H.: 'Generalized free products with amalgamated
subgroups I', Amer. J. Math. 70 (1948),590-625.
[A3] NEUMANN, H.: 'Generalized free products with amalgamated
subgroups II', Amer. J. Math. 71 (1949), 491-540.
AMS 1980 Subject Classification: 08B25, 20AXX
AMICABLE NUMBERS - A pair of natural numbers
each one of which is equal to the sum of the proper
divisors of the other, i.e. of the divisors other than the
number itself. This definition is found already in
Euclid's Elements and in the works of Plato. One pair
of such numbers only - 220 and 284 - was known to
the ancient Greeks; the sums of their divisors are equal,
AMPLE SHEAF
respectively, to
1+2+4+5+10+11+20+22+44+55+110 = 284,
1+2+4+71 + 142 = 220.
L. Euler discovered about 60 pairs of amicable
numbers, while the use of electronic computers yielded
a few hundreds of such numbers. It is not known, how-
ever, if there exists a pair of amicable numbers one of
which is even while the other is odd.
A.!, Galochkin
Editorial comments. For some very large pairs see [A1].
References
[A1] RmLE, H.J.J. TE: 'New very large amicable pairs', in Proc.
Number Theory Noordwijkerhout, 1983, Lect. Notes in Math.,
Vol. 1068, Springer, 1983, pp. 210-215.
AMS 1980 Subject Classification: 10A40
AMPLE SHEAF - A generalization of the concept of
an ample invertible sheaf. Let X be a Noetherian
scheme over a field k, and let C be a locally free sheaf
on X (that is, the sheaf of sections of some algebraic
vector bundle E~X). The sheaf C is called ample if for
each coherent sheaf on X there exists an integer no,
depending on , such that the sheaf S n C for
n ~no is generated by its global sections (here sn C
denotes the n-th symmetric power of C).
A locally free sheaf C on X is ample if and only if
the invertible tautological sheaf !l'( C) on the projectiv-
ization P(E) of the bundle E is ample. Another cri-
terion of ampleness is that for each coherent sheaf
on X there must exist an integer no, depending on ,
such that the cohomology group Hi(X, sn C) is
zero for n~no and i>O. If the sheaves C and are
ample then C is an ample sheaf [1]. If X is a non-
singular projective curve, then a sheaf C on X is ample
if and only if C and all its quotient sheaves have posi-
tive degree [2]. The tangent sheaf on pN is ample for
any N (see [1]). The converse also holds: Any non-
singular N-dimensional algebraic variety with an ample
tangent sheaf is isomorphic to pN (see [1], [3]).
References
[1] HARTSHORNE, R.: 'Ample vector bundles', Publ. Math. IRES
29 (1966),319-350.
[2] HARTSHORNE, R.: 'Ample vector bundles on curves', Nagoya
Math. J. 43 (1971), 73-89.
[3] DEMAZURE, M.: 'Characterisations de l'espace projectif (conjec-
tures de Hartshorne et de Frankel)', in Sem. Bourbaki 1979/80,
Lect. Notes in Mathematics, Vol. 842, Springer, 1981, pp. 11-
1~ V.A. Iskovskikh
Editorial comments. The theorem stated in the last line
of the text is due to S. Mori [A 1].
References
[A1] MORl, S.: 'Positive manifolds with ample tangent bundles',
Ann. of Math. 110 (1979), 593-606.
AMS 1980 Subject Classification: 14F05
165
AMPLE VECTOR BUNDLE
AMPLE VECI'OR BUNDLE - An algebraic or ana-
lytic vector bundle with an ample sheaf of regular
(resp. analytic) sections (see Ample sheaf; Positive vec-
tor bundle).
o.A. Ivanova
AMS 1980 Subject Classification: 14F05
AMPLITUDE OF AN ELLIPTIC INTEGRAL - The
variable 4>, considered as a function of z, in an elliptic
integral of the first kind
z = F(.p,k) = j dt
o VI-k 2 sin2 t
ill the normal Legendre form. The concept of the
amplitude of an elliptic integral and the notation
4>=amz were introduced by e.G.J. Jacobi in 1829. The
amplitude of an elliptic integral is an infinite-valued
periodic function of z. The basic elliptic Jacobi func-
tions sinamz=snz, cosarnz=cnz, 6.amz=dnz are all
single-valued. It is convenient, however (e.g. for pur-
poses of tabulation), to consider an elliptic integral as a
function F(4), k) of the amplitude 4> and the modulus k.
See also Jacobi elliptic functions.
E.D. Solomentsev
AMS 1980 Subject Classification: 33A25
ANALLAGMATIC GEOMETRY - A geometry that
studies the properties of figures which are invariant
with respect to circle transformations (d. Circle
transfonnation) of the extended plane, i.e. of the
Euclidean plane completed by a point at infinite dis-
tance.
A.B. Ivanov
AMS 1980 Subject Classification: 51 M05
ANALYTIC CAPACITY, analytic measure, Ahlfors
analytic measure - A set function in the plane, intro-
duced by L. Ahlfors [1], which is an analogue of the
logarithmic capacity, and which is suited for the charac-
terization of sets of removable singularities of bounded
analytic functions. Let E be a bounded closed set in the
plane, let A (E) be the set of functions which are ana-
lytic outside E, vanish at infinity and are bounded
everywhere outside E by the constant 1. Let
y(E ;f)=limzj(z) as z....,.oo. The number
y(E)=sup{ I y(E;f) I: fEA(E)}, is called the analytic
capacity of the set E. The analytic capacity of an arbi-
trary set is usually defined as the supermum of the ana-
lytic capacities of its bounded closed subsets.
If E is a bounded closed set, then a necessary and
sufficient condition in order that each function which is
analytic and bounded outside E be extendable to E, is
that the analytic capacity of E be zero (the Ahlfors
theorem).
There also exist several concepts related to analytic
166
capacity which are suited. to the metrics of various
other spaces of analytic functions (see, for example, [2]
and [3]).
The concept of analytic capacity proved to be well
suited to certain problems in approximation theory,
where the solution of a number of fundamental prob-
lems is formulated in terms of analytic capacity. Thus,
[4], any continuous function on a bounded closed set E
in the plane can be uniformly approximated by rational
functions to any desired degree of accuracy, if and only
if the equality
y(J~ \ E) = y(J~) = 8
holds for any disc 118 of radius 8.
References
[l] AHLFORS,L.V.: 'Bounded analytic functions', Duke Math. J.
14(1947), I-II.
[2] GARABEDIAN, P.R.: 'The classes 4 and conformal mapping',
Trans. Amer. Math. Soc. 69 (1950), 392-415.
[3] SINANYAN, S.O.: 'Approximation by analytic functions and
polynomials in a real mean', Dokl. Akad. Nauk ArmSSR 35,
no. 3 (1962), 107-112 (in Russian).
[4] VlTUSHKIN, A.G.: 'Analytic capacity of sets in problems of
approximation theory', Russian Math. Surveys 22, no. 6 (1967),
139-200. (Uspekhi Mat. Nauk 22, no. 6 (1967), 141-199)
A. G. Vitushkin
Editorial comments. A good reference to the general
area is [A1].
References
[A 1] GAMELIN, T.W.: Uniform algebras, Prentice-Hall, 1969.
AMS 1980 Subject Classification: 31 A 15, 31 C15
ANALYTIC COMPLEMENT - An obsolete term
denoting a set which is complementary to an s1' -set (on
the number axis).
P.S. Aleksandrov
AMS 1980 Subject Classification: 04A 15
ANALYTIC CONTINUATION of a function - An
extension of a function fa, already defined on a certain
subset E of a complex manifold M, to a function f
which is holomorphic on a certain domain D eM con-
taining E such that the restriction fl E =fo of f to E
coincides with fa. The starting point in the theory of
analytic continuation is the concept of an (analytic) ele-
ment, i.e. a pair (D,f), where D is a domain in M and
f is a holomorphic function on D. One says that the
elements (Do,fo) and (D l,fd are direct analytic con-
tinuations of each other through a connected com-
ponent 6. of the set Do n
D 1 if fa It. =fl It.. By defini-
tion, an element (Do,fo) continues analytically to a
boundary point ~EaD eM if there exists a direct ana-
lytic continuation (D 1 ,ij) of the element (Do,fo)
through 6. such that ~ E 6. n
D 1. A maximal analytic
continuation of CDo,fo) (in /'vf) is an element CD, f)
which continues h analytically to the domain D ~ Do,
 ANALYTIC CONTINUATION

but which cannot analytically be continued to any parameter-dependent integrals such as Cauchy-type
boundary point of D. The maximal analytic continua- integrals (cf. Cauchy integral); the Laplace integral; the
tion of (Do,fo) in M is unique, but does not always Borel transform; change of variable in a power series,
exist. In order to overcome this drawback one intro- special methods of summation of power series (the
duces the concept of a covering domain over M (a Borel expansion into a polynomial series converging in
Riemann surface in the case M = C), which is con- a maximal polygon (cf. Borel summation method), the
structed from the elements that are analytic continua- Mittag-Leffler series converging in a maximal star (cf.
tions of (Do,fo). An element (D, f) is called an ana- Star of a function element; Mittag-Leffler summation
lytic continuation of an element (Do,fo) if there exists a method), etc., the Riemann - Schwarz reflection method
finite chain of elements (Di , J;), i =0, ... , n, and (cf. Riemann - Schwarz principle), functional and dif-
corresponding connected components !1i in Di n Di + I ferential equations satisfied by a function (e.g. the
such that (Dn,fn)=(D,j) and (Di,J;), (Di+I,J;+I) are equation f(z + 1) =z f(z) for the gamma-function, con-
direct analytic continuations of each other through !1i. ditions of periodicity, evenness, symmetry, etc.), and
One says that the holomorphic function fa, initially analytic expressions in terms of known functions.
defined in the domain Do, continues analytically to a The subject of analytic continuation also comprises
point z EM if there exists an analytic continuation studies on the relation between the initial element of an
(D, f) of (Do,fo) such that z ED. One introduces an analytic function (a Taylor series) and the properties of
equivalence relation among the elements which are con- the complete analytic function generated by this element
tinuations of fo to the point z: (D',f}-(D",f") if [1]. Results have been obtained on singular points (cri-
zED' n D" and / =/' in a neighbourhood of z. On teria for singular points, the Hadamard theorem on pro-
the set of equivalence classes Df (for all possible z) ducts, the Fabry theorem on quotients) and on singular
there is a natural way of introducing the topology and curves (theorems on gaps and non-extend ability beyond
complex structure of a covering domain over M. The the boundary of the disc of convergence, e.g. the
function fo is lifted to Df in a natural manner (its value Hadamard theorem on gaps, the Fabry theorem on
on the equivalence class at z containing (Do,fo) is set gaps, etc.), theorems on over-convergence and on rela-
equal to fo(z ; it continues analytically to all of Df tions between analytic continuation of a power series
and, in the sense defined above, it does not continue to and properties of an entire function which defines its
any boundary point of Df over M. coefficients, problems of meromorphic continuation,
If M is the complex plane C or, more generally, the meromorphic continuation by Pade approximation, etc.
complex space en, n ~ 1, this process of analytic con- The field of analytic continuation also includes
tinuation can be described more simply. A canonical theorems on removable singularities (removability of an
element is a pair (Da,fa) where a Ecn and fa is a power . isolated singularity of a bounded holomorphic function,
series centred at the point a with a non-empty domain removability of a rectifiable singular curve under condi-
of convergence Da. A canonical element (Db,fb) is an tions of continuity, etc.), as well as a large class of
analytic continuation of (Da,fa) along a path theorems on the simultaneous continuation of holo-
y: [0, l]~en if there exists a family of canonical ele- morphic functions of several complex variables. The
ments (D,,!t), O~t~l, with centres a,=z(t), such that space en, n> 1, comprises domains for which every
(Do,fo)=(Da,fa), (D1,f1)=(Db,fb), and for any toEY holomorphic function can be extended to a larger
the elements (D,,!t) are direct analytic continuations of domain (this phenomenon is absent in the one-
(D'o,!to) for all t sufficiently near to to. The family dimensional case). It is therefore an important problem
(D,,!t) is in fact uniquely determined. If y" O~'T~ 1, is in the theory of analytic continuation of functions of
a continuous family of paths in Cn with common end several complex variables to describe these larger
points a and b, and if (Da,fa) continues analytically domains - the so-called envelopes of holomorphy. Thus,
along each y" then the result (Db,fb) does not depend there are descriptions of the hulls (envelopes) of holo-
on 'T (the monodromy theorem). In the case of en the morphy for Hartogs domains, n-circular and tube
canonical elements (Da,fa) obtained by analytic con- domains, theorems on the removability of compact
tinuation along all possible paths in en become the singularity sets and of singularity sets of codimension
points of Df ; fa is lifted to Df and continues analyti- ~2, the Bogolyubov theorem on the 'edge-of-the-wedge'
cally throughout Df to a holomorphic function f, while and Vladirnirov's theorem on C-convex hulls [3] (cf.
Df is the domain of holomorphy of f Holomorphic envelope; Hartogs domain; Tube domain).
This general process of analytic continuation is not A number of effective methods for the construction of
very effective in practice, and for this reason many spe- envelopes of holomorphy are available [3].
cial methods of analytic continuation are employed. The problem of analytic continuation of functions of
These include various analytic representations: a real variable can be reduced to that of holomorphic

167
ANALYTIC CONTINUATION
functions, since for any domain G eRn and any func-
tion f an~lytic in G there are a domain D c en and a
~unction f holomorphic in D such that D Rn = G andn
flc=j
References
[I) BIEBERBACH, L.: Analytische Fortsetzung, Springer, 1955.
[2] MARKUSHEVICH, A.I.: Theory offunctions of a complex variable,
1-3, Chelsea, 1977 (translated from the Russian).
[3] VLADIMIROV, V.S.: Methods of the theory offunctions of several
complex variables, M.LT., 1966 (translated from the Russian).
E.M. Chirka
Editorial comments.
References
[A 1] HORMANDER, L.: Introduction to complex analysis in several
variables, v. Nostrand, 1966.
[A2) GRAUERT, H. and FRlTZSCHE, K.: Several complex variables,
Springer, 1976 (translated from the German).
AMS 1980 Subject Classification: 30840, 32A30
ANALYTIC CURVE, analytic arc - A curve K in an
n-dimensional Euclidean space Rn, n;;:;' 2, which has an
analytic parametrization. This means that the coordi-
nates of its points x = (x I, . . . , x n ) can be expressed as
analytic functions of a real parameter x, =x,(t),
i = I, ... , n, a ~ t ~ /3, i.e. in a certain neighbourhood
of each point to, a~to~f3, the functions x,(t) can be
represented as convergent power series in t-to, and
the derivatives x;(t 0), i = I, ... , n, do not simultane-
ously vanish at any point of the segment [a,f3l. This
last condition is sometimes treated separately, and an
analytic curve which satisfies it is called a regular ana-
lytic curve. An analytic curve is called closed if
x;(a)=x;(f3), i=l, ... ,n.
An analytic curve in the plane C=C 1 of the complex
variable z = x I + ix 2 can be represented as a complex-
analytic function of a real parameter z = j(t), a~t~/3,
j'(t)=F0 on [a, 13], If the analytic curve is located in a
domain DeC, then a conformal mapping of D into
any domain will also yield an analytic curve. If the set
of intersection points of two analytic curves is infinite,
these analytic curves coincide.
In general, in a complex space en, n;;:;' I, the complex
coordinates Z; of the points of an analytic curve can be
represented as analytic functions of a real parameter
z,=z,(t), a~t~/3, i=l, ... ,n. It should be noted,
however, that if n> 1, the term 'analytic curve' may
sometimes denote an analytic surface of complex
dimension one.
On a Riemann surface S an analytic curve K can be
represented as j(t)=If;(<P(t)), where Z=If;(P) is a local
uniformizing parameter of the points P on Sand j (t)
is an analytic function of a real parameter in a neigh-
bourhood of any point to E[a, /3].
References
[1] MARKUSHEVICH, A. I.: Theory of/unctions of a complex variah/e.
1. Chelsea, 1977 (translated from the Russian).
168
[2) SHABAT, B.V.: Introduction to complex analySiS, Moscow, 1969,
Chapt. 3 (in Russian).
Editorial comments.
E.D. Solomentsev
References
[A1] GUNNING, R.C. and ROSSI, H.: Analytic functions of several
complex variables, Prentice-Hall, 1965.
AMS 1980 Subject Classification: 53A04, 32C25
ANALYTIC DIFFERENTIAL - See Differential on a
Riemann surface.
AMS 1980 Subject Classification: 14H05, 30F30
ANALYTIC EXPRESSION (formula) - The totality of
operations to be performed in a certain sequence on the
value of an argument and on the constants in order to
obtain the value of the function. Every function in one
unknown x with not more than a countable number of
discontinuities has an analytic expression A (x) involv-
ing only three operations (addition, multiplication,
passing to the limit by rational numbers), performed
not more than a countable number of times, starting
from an argument x and from the constants, e.g.
_ 00 n x2n+1
SIllX - n~o( -I) (2n + I)!'
If there is at least one analytic expression describing a
given function, there are infinitely many such expres-
sions. Thus, the function which is identically equal to
zero is expressed by the series
0=
00
~xx-n+1
n I( )
n=1 n!
and from any analytic expression A (x) it is always pos-
sible to obtain another one which is identically equal to
the first:
A(x)+B(x) [n~1 xn-I~~-n) +11,
where B (x) is an arbitrary analytic expression.
References
[I] LUZIN, N.N.: Theory of functions 0/ a real variahle, Moscow,
1948 (in Russian).
B. V Kutuzov
AMS 1980 Subject Classification: 26A09
ANALYTIC FUNCfION - A function that can be
locally represented by power series. The exceptional
importance of the class of analytic functions is due to
the following reasons. First, the class is sufficiently
large; it includes the majority of functions which are
encountered in the principal problems of mathematics
and its applications to science and technology.
Secondly, the class of analytic functions is closed with
respect to the fundamental operations of arithmetic,
algebra and analysis. Finally, an important property of
 ANAL YTIC FUNCTION

an analytic function is its uniqueness: Each analytic sense of complex analysis or in the sense of C) at z;
function is an 'organically connected whole', which A =/ (z) is the derivative of f at z, and
represents a 'unique' function throughout its natural Afj,z = j'(z)dz = df(z)
domain of existence. This property, which in the 18-th
century was considered as inseparable from the very is its differential at that point. A function f which is
notion of a function, became of fundamental signifi- differentiable at every point of D is called differentiable
cance after a function had come to be regarded, in the in the domain D.
first half of the 19-th century, as an arbitrary One may compare the concepts of differentiability of
correspondence. The theory of analytic functions ori- f considered as a function of two variables (in the sense
ginated in the 19-th century, mainly due to the work of of R2) and in the sense of C. In the former case the dif-
A.L. Cauchy, B. Riemann and K. Weierstrass. The ferential dfhas the form
'transition to the complex domain' had a decisive effect 11
ax dx+ 11
ay dy,
on this theory. The theory of analytic functions was
constructed as the theory of functions of a complex where
variable; at present (the 1970's) the theory of analytic J!...
ax = ~+i~ lL
ax ax 'ay = ~+i~
ay ay ,
functions forms the main subject of the general theory are the partial derivatives of f Passing from the
of functions of a complex variable. independent variables x, y to the variables z, Z, which
There are different approaches to the concept of may formally be considered as new independent vari-
analyticity. One definition, which was originally pro- ables, related to the old ones by the equations
posed by Cauchy, and was considerably advanced by z=x+ry, z=x-ry (from this point of view, the func-
Riemann, is based on a structural property of the func- tion f may also be written as f (z, Z)) and expressing dx
tion - the existence of a derivative with respect to the
and dy in terms of dz and di according to the usual
complex variable, i.e. its complex differentiability. This
rules of differential calculus, one can write df in its
approach is closely connected with geometric ideas.
complex form:
Another approach, which was systematically developed
by Weierstrass, is based on the possibility of represent- dif = Jf.. dz + lL di
ing functions by power series; it is thus connected with az az'
the analytic apparatus by means of which a function where
can be expressed. A basic fact of the theory of analytic
functions is the identity of the corresponding classes of Jf.. = 1. [J!... -
az 2 ax ilL] lLaz = 1. [J!...ax + ilL]
ay' 2 ay'
functions in an arbitrary domain of the complex plane.
Exact definitions are given below. Let D be a domain are the (formal) derivatives of f with respect to Z and Z,
in the complex plane C. If to each point Z ED there has respectively. It is seen, accordingly, that f is differenti-
able in the sense of C if and only if it is differentiable in
been assigned some complex number w, then one says
the sense of R2 and if the equation af / az=o is satis-
that on D a (single-valued) function f of the complex
variable Z has been defined and one writes: w = f (z), fied, which in expanded form may be written as
ZED (or f:D~C). The function w=f(z)=f(x+iy) ~-~ ~--~
may be regarded as a complex function of two real ax - ay' ay - ax .
variables x and y, defined in the domain D CR2 (where If f is differentiable in the sense of C in D, the latter
R2 is the Euclidean plane). To define such a function relations are satisfied at all point of the domain; they
is tantamount to defining two real functions are called the Cauchy- Riemann equations. These equa-
tions occurred already in the 18-th century in J.L.
U = 4>(x,y), v = l/J(x,y), {x,y)ED (w = u+iv).
d'Alembert's and L. Euler's studies on functions of a
Having fixed a point Z ED, one gives Z the increment complex variable. The initial definition may now be
Az =.ix + i Ay (such that Z + Az E D) and considers the rendered more precise as follows. A function f, defined
corresponding increment of the function f in a domain D, is said to be holomorphic (analytic) at a
!!.f{z) = fez + fj,z) - f(z). point Zo ED if there exists a neighbourhood of this
point in which f may be represented by a power series:
If
!!.f{z) = A fj,z + o(fj,z)
fez) = ao +a1 (z - Zo)+ ... +an{z - ZoY + ....
as Az ~o, or in other words, if If this requirement is satisfied at every point Z 0 of D,
lim MEl =A the function f is said to be holomorphic (analytic) in the
&~o fj,z domain D.
exists, the function f is said to be differentiable (in the A function f which is holomorphic at a point Zo ED

169
ANALYTIC FUNCTION

is differentiable at that point. In addition, the sum of a formula, which expresses the values of an analytic func-
convergent power series has derivatives of all orders (is tion inside a domain in terms of its values on the boun-
infinitely differentiable) with respect to the complex
variable z; the coefficients of the series may be
dary:
1
fez) - -2' J- 1niill!R
,zED.
7TI D ~-z
expressed in terms of the derivatives of f at Zo by the
formulas an = j<n)(zo) / n!. The power series, written in Here, D is a domain whose boundary aD consists of a
the form finite number of non-intersecting rectifiable curves (the
l(zo) j<n)(zo);, orientation of aD is assumed to be positive with respect
f(zo)+-l-! -(z-zo)+ ... + n! (z-zo)+"', to D), and L. is a function which is analytic in some
domain G-::JD=D UaD. This formula makes it possi-
is known as the Taylor series of f at zoo Thus, holomor-
ble, in particular, to reduce the study of many prob-
phy of a function f in a domain D means that it is
lems connected with analytic functions to the
infinitely differentiable at any point in D and that its
corresponding problems for a very simple function -
Taylor series converges to it in some neighbourhood of
the Cauchy kernell/(t-z), tEaD, ZED. For more
this point.
details see Integral representation of an analytic func-
On the other hand, the following noteworthy fact is
tion.
established in the theory of analytic functions: If a
A very important property of analytic functions is
function f is differentiable in a domain D, it is holo-
expressed by the following uniqueness theorem: Two
morphic in this domain (for a single point this state-
functions which are analytic in a domain D and which
ment is not true: f (z) = 1 Z 12 = zz is differentiable at
coincide on some set with a limit point in D, coincide
Zo =0, but is nowhere holomorphic). Thus, the concepts
throu.ghout D (are identical). In particular, an analytic
of complex differentiability and holomorphy of a func-
functIOn f (z), z ED, which is not identically zero can
tion in a domain are identical; each one of the follow-
only have isolated zeros in D. If, in addition, z 0 is a
ing properties of a function f in a domain D - dif-
zero of f, then one has, in some neighbourhood U(zo)
ferentiability in the sense of C, differentiability in the
of zo, f(z)=(z-zo)"g(z), where p~l is a natural
sense of R2 together with satisfaction of the
number (called the multiplicity of the zero of f at zo),
Cauchy - Riemann equations, holomorphy - may
while g(z) is a analytic function in U(zo).
serve as definition of analyticity of f in this domain.
An important role in the theory of analytic functions
One other characteristic of an analytic function is
is played by the points at which the function is not
connected with the notion of an integral. The integral
analytic - the so-called singular points of the analytic
of a function f=cf>+il[; along an (oriented rectifiable)
function. Here, only isolated singular points of (single-
curve f: z=z(t), tE[a,p), may be defined by the for-
valued) analytic functions are considered; for more
mula:
f3 details cf. Singular point. If f is an analytic function in
[fez )dz = [fIz(t)]z' (t) dt an annulus of the form 0< 1 z -Zo 1 <p, it may be
expanded there in a Laurent series
or by means of a curvilinear integral: +oc

[f(z)dz = [[cf>dx -""dy]+i [[""dx +cf>dy].

fez) = ~ an{z - Zo)",
n= -00

which contains, as a rule, not only positive but also

A key result in the theory of analytic functions is
negative powers of z - zoo If there are no terms with
Cauchy's integral theorem: If f is an analytic function in
negative powers in the series (an =0 for
a domain D, then Jrf (z) dz = 0 for any closed curve
n = -1, -2, ... ), Zo is called a regular point of f (a
reD bounding a domain inside D. The converse
removable singular point). At a regular point there also
result (Morera's theorem) is also true: If f is continuous
exists a finite limit
in a domain D and if frf (z) dz = 0 for any such curve
r, then f is an analytic function in D. In particular, in limf(z)
.?~:o
= ao;
a simply-connected domain, those and only those con-
tinuous functions f are analytic, whose integral along Putting f(zo)=aQ, one obtains an analytic function in
any closed curve reD is zero (or. which is the same the whole disc z - Zo <po If the Laurent series of the
1 1

function contains only a finite number of terms with

thing, the integral along any curve r connecting two
negative powers of z - z 0:
arbitrary points z 1,22 ED does depend only on the
points z I and Z2 themselves and not on the shape of
00

fez) = ~an(z -zo)". /.I. < 0, a~ =1= 0.

the curve). This characterization of analy'tic functions
forms the basis of many of their applications. the point Zo is called a pole of f (of multiplicity fL); a
Cauchy's integral theorem yields Cauchy's integral pole z 0 is characterized by

170
 ANAL YTIC FUNCTION

limf(z) = 00. The simplest class of analytic functions consists of

Z---+-Zo
the functions which are holomorphic in the whole
The function f has a pole of multiplicity j.t at the point
plane; such functions are called entire functions. Entire
Z 0 if and only if the function 1 / f has a zero of multi-
functions are represented by series
plicity j.t at this point. If the Laurent series contains an ao +alz + ... +anz n + ... which are convergent in
infinite number of negative powers of Z - Zo (an:;;t=O for
the whole plane. This class includes the polynomials in
an infinite set of negative indices n), then Zo is called
z, the functions
an essential singular point; at such points there is no
finite and no infinite limit for f The coefficient a ~ I in eZ = 1+-+-+
Z
...
Z2 zn
+ - + ...
the Laurent series for f with centre at the isolated I! 2! n!
singular point Z 0 is called the residue of f at z 0 : . Z3 z2n+1
SlllZ = z--+ ... +(-1)" + ...
3! (2n+l)!
a_I = res[f(z);z=zoJ. z2n
It can be defined by the formula
cosz = Z2
1 - - + ... +(-1)"--+ ...
2! (2n)!'

res[f(z);z=zol = ~ jf(z)dz, etc.

7Tl r
Weierstrass' theorem states that, for any sequence of
where y={z: I Z-Zo I =p} and p>o is sufficiently complex numbers an, n = 1,2, ... , without limit points
small (so that the disc I Z-Zo I,,;;;p does not contain in C, there exists an entire function F that vanishes at
singular points of f other than zo). The important role the points an and only at these points (among the an
of residues is made clear by the following theorem: If f there may be coincident points, to which correspond
is an analytic function in a domain G, except for some zeros of F of corresponding multiplicity). Here, the
set of isolated singular points, if reG is a contour function F may be represented as a (generally infinite)
bounding a domain D C G and not passing through any product of entire functions each one of which has only
singular points of J, and if z I, . . . ,Zn are all the singu- one zero.
lar points of f inside D, then Functions that are meromorphic in the plane (i.e.
that may be represented as quotients of entire func-
/f(z )dz = 2'7Tik~/es[f(z); z =zd tions) are called meromorphic functions. These include
rational functions, tgz=sinz /cosz,
This theorem is an effective tool in calculating
ctgz =cosz / sinz, elliptic functions, etc.
integrals. See also Residue of an analytic function.
According to Mittag-Lejjler's theorem, for any
The sum of the terms of the Laurent series for f at
sequence fin EC, n = 1,2, ... , without limit points in
z 0 corresponding to the negative indices n,
C, there exists a meromorphic function G with poles at
-I
~ an(z -zo)", the points fin and only at those points, such that its
n= -00
principal parts at the points fin coincide with pre-
is known as the principal part of the Laurent series (or assigned polynomials in 1/ (z - fin). The function G
of the function 1> at the point zo0 This principal part may be represented as a (usually infinite) sum of mero-
determines the nature of the singularity of f at z o. morphic functions, each one with a pole at a single
Functions which are representable as a quotient of point only.
two functions that are holomorphic in a domain Dare Theorems on the existence of a holomorphic function
called meromorphic in the domain D. A function which with pre-assigned zeros and of meromorphic functions
is meromorphic in a domain is holomorphic in that with pre-assigned poles and principal parts are also
domain, except possibly at a finite or countable set of valid for an arbitrary domain DeC.
poles; at the poles the values of a merom orphic func- In the study of analytic functions the related
tion are considered to be infinite. If such values are geometric notions are also of importance. If f: D~C is
allowed, then meromorphic functions in a domain D an analytic function, the image feD) of the domain D
may be defined as functions that in a neighbourhood of is also a domain (principle of preservation of domains).
each point zoED can be represented by a Laurent If / (z 0 ):;;t=0, the mapping f preserves angles at z 0 both
series in z - z 0 with a finite number of terms involving in value and in sign, i.e. it is conformal. Thus, there
negative powers of Z - z 0 (depending on z 0) in a neigh- exists a close connection between analyticity and the
bourhood of each point Zo ED. important geometric notion of conformal mapping. If f
Both holomorphic and merom orphic functions in a is an analytic function in D and f(z'):;;t=f(z") for
domain D are often designated as analytic in the domain z' :;;t=z" (such functions are called univalent), then
D. In this a case holomorphic functions are also said to / (z):;;t=O in D and f defines a one-to-one and conformal
be regular analytic or simply regular functions. mapping of the domain D onto the domain feD).

171
ANALYTIC FUNCTION
Riemann's theorem, which is the fundamental theorem
in the theory of conformal mappings, says that on any
simply-connected domain whose boundary contains
more than one point there exist univalent analytic func-
tions which conformally map this domain onto a disc
or a half-plane (cf. Conformal mapping; Univalent func-
tion).
The real and imaginary parts of a function f = </> + i1P
which is holomorphic in a domain D satisfy the
Laplace equation in that domain:
ft
2+
f t -- 0, f i+f i - 0,
ax 2
2 -
ax ay2 ay
i.e. they are harmonic functions (cf. Harmonic func-
tion). Two harmonic functions which are connected by
the Cauchy - Riemann equations are called conjugate.
In a simply-connected domain any harmonic function </>
has a conjugate function 0/ and is thus the real part of
some holomorphic function finD.
The connections with conformal mappings and har-
monic functions form the basis of many applications of
the theory of analytic functions.
A functionf(z), ZEE (ECC being an arbitrary set),
is called analytic at a point ZOE E if there exists a neigh-
bourhood of this point such that f may be represented
by a convergent power series on the intersection of this
neighbourhood with E. The function f is called analytic
on the set E if it is analytic on some open set which
contains E (or, more exactly, if there exist both an
open set containing E and an analytic function F on
this set which coincides with f on E). For open sets the
notion to analyticity coincides with the notion of dif-
ferentiability with respect to the set. However, this is
not the case in general; in particular, on the real line
there exist functions which not only have a derivative,
but which are infinitely differentiable at every point
and are not analytic even at a single point of this line.
The property of connectedness of the set E is necessary
in order that the uniqueness theorem for analytic func-
tions holds. This is why analytic functions are usually
considered in domains, i.e. on connected open sets.
All the preceding refers to Single-valued analytic func-
tions f, considered in a given domain D (or on a given
set E) of the complex plane. In considering the possible
extension of a function f, as an analytic function, to a
larger domain, one arrives at the concept of the ana-
lytic function considered as a whole, i.e. throughout its
whole natural domain of existence. If the function is
thus extended, its domain of analyticity becomes larger,
and may overlap itself, supplying new values of the
function at points in the plane where it already was
defined. Accordingly, an analytic function considered
as a whole is generally multi-valued. Many problems in
analysis (inversion of a function, the determination of a
172
primitive and the construction of an analytic function
with a given real part in multiply-connected domains
(cf. Multiply-connected domain), the solution of alge-
braic equations with analytic coefficients, etc.) require
the study of multi-valued functions; such functions
include Z I / n, In z, arctg z, arcsin z, algebraic functions,
etc. (cf. Algebraic function).
A regular process which yields the complete analytic
function, considered throughout its natural domain of
existence, was proposed by Weierstrass; it is known as
Weierstrass' method of analytic continuation.
The initial concept is that of an element of an ana-
lytic function, viz. a power series with a non-zero
radius of convergence. Such an element Wo:
aO+al(z-zO)+ ... +an(z-zo)"+ ... ,
defines a certain analytic function f on its disc of con-
vergence Ko. Let z\ be a point of Ko different from zo0
Expanding f in a series with centre at z I, one obtains a
new element WI:
bo+bl(z-Zl)+ ... +bn(Z-ZI)n+ ... ,
whose disc of convergence will be denoted by K I. In
the intersection of Ko and KI the series WI converges
to the same function as the series W o. If K I extends
beyond the boundary of Ko, the series WI defines the
function determined by Wo on some set outside Ko
(where the series Wo is divergent). In such a case the
element W \ is called a direct analytic continuation of
the element W 0 0 Let W o, . .. , WN be a chain of ele-
ments in which Wn + I is a direct analytic continuation
of Wn (n =0, ... , N -1); the element W N is then said
to be an analytic continuation of the element Wo (by
means of the given chain of elements). When the centre
of the disc KN belongs to Ko it may happen that the
element WN is not a direct analytic continuation of the
element W 0 0 In such a case the sums of the series Wo
and WN will have different values in the intersection of
Ko and K N; thus analytic continuation may lead to
new values of the function inside K o.
The totality of all elements which may be obtained
by analytic continuation of an element Wo forms the
complete analytic function (in the sense of Weierstrass)
generated by Wo; the union of their discs of
convergence represents the (Weierstrass) domain of
existence of this function. It follows from the unique-
ness theorem for analytic functions that an analytic
function in the sense of Weierstrass is completely deter-
mined by the given element W 0 The initial element
0
may be any other element belonging to this function;
the complete analytic function will not be affected.
A complete analytic function f, considered as a func-
tion of the points in the plane belonging to its domain
of existence D, is generally multi-valued. In order to
eliminate this feature, f is considered not as a function

of the points in the plane domain D, but rather as a
function of the points on some multi-sheeted surface R
(lying above D) such that to each point of D there
correspond as many points of the surface R (projecting
onto the given point of D) as there are different ele-
ments with centre at this point in the complete analytic
function f; on the surface R the function f becomes
single-valued. The idea of passing to such surfaces is
due to Riemann, and the surfaces themselves are
known as Riemann surfaces. The abstract definition of
the notion of a Riemann surface has made it possible
to replace the theory of multi-valued analytic functions
by the theory of single-valued analytic functions on
Riemann surfaces (cf. Riemann surface).
Now fix a domain d belonging to the domain of
existence D of the complete analytic function f, and fix
some element W of f with centre at a point in d. The
totality of all elements which may be obtained by ana-
lytic continuation of W by means of chains with cen-
tres belonging to d is called a branch of the analytic
function f A branch of a multi-valued analytic function
may tum out to be a single-valued analytic function in
the domain d. Thus, arbitrary branches of the functions
Z I / n and In Z which correspond to an arbitrary simply-
connected domain not containing the point 0, are
single-valued functions. The function Z I / n has exactly
n different branches in such a domain, while Inz has an
infinite set of such branches. The selection of single-
valued branches (using some cuts in the domain of
existence) and their study by the theory of single-
valued analytic functions constitute one of the princi-
pal methods of studying specific multi-valued analytic
functions.
A.A. Gonchar
Analytic functions of several complex variables. The
complex space en, conslstmg of the points
Z=(Zb ... ,zn), Zk=Xk+ryb is a vector space over the
field of complex numbers with the Euclidean metric
Iz =
1 [k~IIZk 12)1/2, C = C.
It differs from the 2n-dimensional Euclidean space R2n
by a certain asymmetry: on passing from R2n to en (i.e.
on introducing a complex structure in R2n), the coordi-
nates are subdivided into pairs which appear in the
complex combinations Zk =Xk + iYk.
If a complex function f=cp+j,f is defined in a
domain D ce n and is differentiable at all points ZED
in the sense of R2n (i.e. as a function of the 2n real
variables Xk and Yk), its differential may be represented
in the form
ANALYTIC FUNCTION
bois af / aZk and af / aZk are defined as in the case of
the plane. If df is of the form
df = ~
nxdZ b
k=1 aZk
i.e. is a complex linear function in dz \, ... , dz n , the
function f is said to differentiable in the sense of en or
holomorphic or analytic in the domain D.
Thus, the condition of holomorphy of f in a domain
D c en consists of the condition that it be differentiable
in the sense of R2n and that it satisfy the system of n
complex equations af /azk=o (k=l, ... ,n), which is
equivalent to the system of 2n first-order partial dif-
ferential equations
aXk --a'
~ Jf.. ~
a --
!Yk!Yk
_~ - ... ,n)
a (k-l,
Xk
(Cauchy- Riemann system).
In the case of space (n > 1) as distinct from that of
the plane (n = 1) this system is overdetermined, since
the number of equations is larger than that of the
unknown functions. It remains overdetermined on pass-
ing to the (geometrically more natural) spatial analogue
of a holomorphic function of one complex variable - a
holomorphic mapping f, which is realized by a system
f =ifl, . . . ,In) of n functions fk which are holo-
morphic in a domain D C en . The mapping
f: D~G =f (D) is called biholomorphic if it is one-to-
one and if it is holomorphic together with its inverse
f- I : G~D. The conditions for holomorphy of a map-
ping f: D~en are expressed by a system of 2n 2 real
equations involving 2n real functions. The overdeter-
minacy of the conditions of holomorphy for n > 1 is the
cause of a number of effects typical of the spatial case
- such as the absence of a spatial analogue of the
Riemann theorem on the existence of conformal map-
pings. According to Riemann's theorem, if n = 1, any
two simply-connected domains whose boundaries do
not reduce to a single point are isomorphic. However, if
n > 1, even such simple simply-connected domains as
the ball {I Z 1< I} and the product of discs (polydisc)
{ I z. I < 1: p= 1, ... , n} are non-isomorphic. The
non-isomorphism is brought to light on comparing the
groups of automorphisms of these domains (i.e. their
biholomorphic mappings onto themselves, cf. Biholo-
morphic mapping) - the groups prove to be algebrai-
cally non-isomorphic, whereas a biho10morphic map-
ping of one domain onto another, if it existed, would
establish an isomorphism of these groups. Owing to
this, the theory of biholomorphic mappings of domains
in complex space is substantially different from the
theory of conformal mappings in the plane.
A function lis called holomorphic at a point a E en if
it is holomorphic in some neighbourhood of this point.
According to the Cauchy - Riemann criterion, a func-
173
ANALYTIC FUNCTION

tion of several variables which is holomorphic at a
point a is holomorphic with respect to each variable (if
the values of the other variables are fixed). The con-
verse proposition is also true: If, in a neighbourhood of
some point, a function f is holomorphic with respect to
each variable separately, then it is holomorphic at this
point (Hartogs' fundamental theorem).
In analogy with the case of the plane the holomor-
phy of a function f at a point a = (a I, . . . ,an) is
equivalent to its expandability in a multiple power series
in a neighbourhood of this point
I(z) = ~Ckldz,-a,/l "'(zn-an/",
k,;;'O
may be locally described as sets of common zeros of
polynomials in the variable Zn, with coefficients from
the ring of holomorphic functions in the other variables
'z. This circumstance permits extensive use of algebraic
methods in the local study of analytic sets.
Cauchy's integral theorem must also be slightly
modified in the spatial case, and is then known as the
Cauchy-Poincare theorem: Let a function f be holo-
morphic in a domain D cen; then, for any (n + 1)-
dimensional surface G compactly imbedded in D, with
piecewise-smooth boundary aG,
/!dz = o.

or, in abbreviated notation, As in the planar case, this integral is defined by a

parametric representation of the given set: If aG has
I(z) = ~ ck(z-ai, the equation Z = Z (t), where the parameter
jk j;;.o
t=(tl, ... ,tn } varies over an n-dimensional cell
where k=(k l , ,kn } is a multi-index of integers
r={ak"$:;;t~f3k: k=I, ... ,n}, then, by definition,
k.-;;:.O, I k I =kl + ... +kn and
(z-ai = (z,-a,)k l (zn-an)k". lI
G
dz = f.
"
a(z" ... ,zn)
JIz(t)] a(
t" ... ,tn
) dt, ... dtn

A holomorphic function is infinitely differentiable, and
the above series is its Taylor series, i.e.
I akl + ... H" [
=
az7' ... aZ~"'
Ck
k,!'" kn !
the derivatives being taken at the point a.
The fundamental facts of the theory of holomorphic
functions of one variable extend to holomorphic func-
tions of several variables, sometimes in an altered form.
An instance of this is the Weierstrass preparation
theorem (d. Weierstrass theorem), which extends the
property of holomorphic functions of one variable to
become zero as an integral power of z -a to the spatial
case. The theorem is formulated as follows: If a func-
tion /=1=0 which is holomorphic at a point a Een is
equal to zero at that point, then it may be represented,
in a certain neighbourhood U (a) (possibly after a non-
degenerate linear transformation of the independent
variables) in the form
I(z) = {(zn -ant +c,('z)(zn -an)k-' + ...
where k;;.l is an integer, c v (1' = I, ... ,k) are func-
tions of 'z = (z j, . . . ,Zn -I) which are holomorphic in a
neighbourhood of the point
'a = (al,'" ,an-dEen-1 (a 'prime' preceding a
letter denotes the projection on the space of the first
n -1 coordinates) and which are equal to zero at a,
while <p is holomorphic and zero-free in U(a).
This theorem is of fundamental importance in the
study of analytic sets (cf. Analytic set), which are
described locally, in a neighbourhood of each one of
The difference between the spatial and the planar cases
consists in the fact that in the former case the dimen-
sion of the surface G is less than that of the domain D,
while in the planar case the dimensions are the same
(n + I =2n).
The spatial analogue of the Cauchy integral formula
can be written in a particularly simple form for poly-
cylinder domains, i.e. for products of plane domains.
Let D=DI X ... XDn be a domain in which Dk is a
domain in the complex plane with a piecewise-smooth
boundary aDk (k = I, ... ,n), while the function f is
holomorphic in a domain which compactly contains D.
Successive application of Cauchy'S formula for one
variable then yields, for any point zED,
I(z) f
= (2~i)" [f~2:~,
where r = aD I X ... X aDn is an n-dimensional surface
in the boundary aD, f=(fl, ... ,fn), and
d~ _ d~, ... d~n
+Ck('Z)}CP(Z),
~-z - (~-z,) (~n-zn)'
However, polycylinder domains are only a very special
class, and in general domains such a separation of vari-
ables is not possible. The role of Cauchy's integral for
arbitrary domains D C en with a piecewise-smooth
boundary is played by the Martinelli - Bochner integral
formula: For any f~ction f which is holomorphic in a
domain containing D, and for any point zED,
fez)
.
= (n -I)!
(27Ti)n
Jim 8(t-z)~~.
I ~-z I ~n
elf)

their points, as sets of common zeros of a certain where df=d~1 ... dL, and
number of functions which are holomorphic at this n

point. By Weierstrass' preparation theorem such sets 8(w) = ~(- I)' t 'w, dWI . dw, 1 dW"1 ... dw".
, 1

174

This is Green's formula for a pair of functions, one of
which is holomorphic in D, while the other is a funda-
mental solution of the Laplace equation in the space
R2n with singular point r=z. If n = 1, this is the ordi-
nary Cauchy integral. If n> 1, the formula differs from
Cauchy's multiple integral for a product of plane
domains in that, first, the integration is not over an n-
dimensional part of the boundary, but over the whole
(2n - 1)-dimensional boundary of the domain, and,
secondly, its kernel (the factor multiplying f under the
integral sign) does not depend analytically on the
parameter z. An analytic kernel, however, is essential in
a number of problems, and it is therefore desirable to
construct an integral formula with such a kernel for as
large a class of domains as possible. An ample supply
of integral formulas, including formulas with an ana-
lytic kernel for many domains, is contained in the gen-
eral Leray formula. This formula is
fez) = (n -l)! (fm 8[wmld~
(27Ti)" tD <~-z, wm>n '
where W=(W1' . . . ,Wn ) is a smooth vector function
depending also on z, dr
and 8 are defined above, and
<r-z,w>=:LZ=l(rk-zdwk; it is assumed that
<r-z,wm>=I=O for any fixed ZED and r running
over aD. The value of the integral in this formula does
not depend on the choice of the vector function wen
(provided that_<r-z,wm>=I=O for all ZED, rEaD),
and if wm = r- z,
this integral coincides with the
Martinelli - Bochner integral. By varying the choice of
w for different classes of domains, the Leray formula
will yield various integral formulas. In the theory of
analytic functions of several variables other integral
representations, which are valid only for certain classes
of domains, are also considered. An important class of
this kind consists of the so-called Wei! domains, which
are a generalization of the product of plane domains.
For such domains one has the Bergman - Weil
representation with a kernel which also depends analyt-
ically on the parameter.
As in the planar case, the study of the singularities of
analytic functions is of fundamental interest; the main
difference between the planar and the spatial cases is
expressed by the Osgood - Brown theorem on the
removability of compact singularities, according to
which any function f which is holomorphic in D \ K,
where D is a domain in cn (n > 1) and K is compact
subset of D which does not subdivide D, extends holo-
morphically to the whole domain D. By this theorem,
holomorphic functions of several variables cannot have
isolated singular points. These are replaced in en
(n > I) by singular sets which are analytic if their
dimension is lower than 2n - 1.
This fact is essential in the theory of multi-
ANAL YTIC FUNCTION
dimensional residues. This theory deals with the prob-
lem of computing the integral of a function f, which is
holomorphic everywhere in a domain D C en except for
an analytic set M, over a closed n-dimensional surface
(J not intersecting M. Since the dimension of the singu-
lar set M is lower than the dimension of D by at least
two, M does not subdivide D. If the surface (J is not
linked with M, i.e. bounds an (n + I)-dimensional sur-
face G compactly belonging to D \ M, then by the
Cauchy - Poincare theorem f"f dz = O. In order to cal-
culate this integral in the general case, it is necessary to
clarify how (J is linked with the singular set M, and to
calculate the integrals over special n-dimensional sur-
faces associated with separate portions of the set M
(residues).
The solution of this problem involves considerable
topological and analytic difficulties. These may often be
overcome by the methods proposed by E. Martinelli
and J. Leray. The Martinelli method is based on the
use of the topological Aleksander- Pontryagin duality
principle, and reduces the study of the n-dimensional
homologies of the set D \ M to the study of the
(n -I)-dimensional homologies of the singular set M.
The Leray method is more general: it is based on the
examination of special homology classes and on the
calculation of certain differential forms (residue forms).
The multi-dimensional theory of residues has also
found applications in theoretical physics (cf. Feynman
integral).
The Osgood - Brown theorem reveals an important
fundamental difference between the spatial and the
planar theories. In the plane one can, for any domain
D, construct a function f which is holomorphic in D
but which cannot be extended analytically beyond its
boundary, i.e. D is the natural domain of existence. In
space the situation is different: thus, the spherical shell
{ 1 / 2 < I z I < I} cannot be the domain of existence of
any holomorphic function, since by the
Osgood - Brown theorem, any function which is holo-
morphic in it will certainly extend analytically to the
entire ball { I z I < I}.
Thus arises the problem of the characterization of the
natural domains of existence of holomorphic functions
- the so-called domains of holomorphy. A simple suffi-
cient condition may be formulated with the aid of the
concept of a barrier at a boundary point of the domain,
i.e. a function fl(z) which is holomorphic in this
domain and which increases without limit as z tends to
r D C en will be a domain of holomorphy if it is possi-
ble to construct a barrier for an everywhere-dense set
of points in its boundary. This condition is satisfied, in
particular, by any convex domain; for any point rEaD
it is sufficient to select, in the (2n - 1)- dimensional
supporting plane to D at the point r, a (2n - 2)-
175
ANAL YTIC FUNCTION

dimensional plane of the form pseudo-convexity of a domain. .

n As regards domains D which are not dom~s. of
L\(z) = ~ak(Zk -Kk) = 0; holomorphy, there arises the problem of desc~bmg
k=l
their envelope of holomorphy, i.e. the smallest domam of
the function h = I / L r will then be a barrier. Conse- holomorphy to which any function holomorphic in D
quently, every convex domain in cn is a domain of
extends analytically. For domains of the simplest types
holomorphy. However, convexity is not a necessary
envelopes of holomorphy can be effectively constru~te?,
condition for holomorphy: A product of plane domains
but in the general case the problem is unsolvable W1t~
is always a domain of holomorphy, and such a product
the class of single-sheeted domains. Under analytIc
need not be convex. Nevertheless, if the notion of con-
continuation of functions beyond the boundary of a
vexity is suitably generalized, it is possible to arrive ~t
given domain D, multi-valuedness may result, which
a necessary and sufficient condition. One such generali-
can be avoided by introducing multi-sheeted covering
zation is based on the observation that the convex hull
domains over D, analogous to Riemann surfaces (cf.
of a set K eRn may be described as the set of points at
Riemann surface). In the class of covering domains, the
which the value of any linear function does not exceed
problem of constructing envelopes of holo~or:rhy .is
the supremum of the values of this function on K. By
always solvable. This problem also has applicatIOns m
analogy, the holomorphically convex hull of a set
theoretical physics, to wit in quantum field theory.
KeD e en is defined by
The transition from the plane to a complex space

Kc: = {ZED: If(z) I ~~~F Ifm 1,fE<'?(D+
where @ (D) denotes the set of all functions which are
holomorphic in D. A domain DeC" is called holo-
morphicallJ: convex if, for every compact subset K of D,
the hull K ('J also is a compact subset of D. Holo-
morphic convexity is a necessary and sufficient ~ond~
tion for a domain of holomorphy. However, this cn-
terion is not very effective, since holomorphic convexity
is difficult to verify.
Another generalization is connected with the notion
of a plurisubhannonic function, which is the complex
analogue of a convex function. A convex function in a
domain D of Rn may be defined as a function for
which the restrictions to the segments in D of the
straight lines x=xo+wt (where xO,wERn and t is a
real parameter) are convex functions of t. ~ real fun~
tion cf>, defined and upper semi-continuous m a domam
D e en, is said to be plurisubharmonic in D if for each
complex line Z = zO + wt (zo, wEC", t is .a .co~plex
parameter), its restriction to the parts of th!s line I~ D
is a sub harmonic function of r
If cf> is tWIce contmu-
ously differentiable, then the condition of 'pluris~b~ar
monicity, in accordance with the rules of dIfferentIatIOn
of composite functions, is that the Hermitian form
n ~_
H(w, w) = ~ _ WjWk
j,k = 1 azJaZ k
- the so-called Levi form - be non-negative.
A domain D e en is called pseudo-convex if the func-
tion In d(z, aD), where d(z, aD) denotes the Euclidean
distance from the point z ED to the boundary aD, is
plurisubharmonic in this domain. Pseudo-convexity .is
also a necessary and sufficient condition for a domam
to be a domain of holomorphy.
In some cases it is possible to verify effectively the
substantially increases the variety of geometrical prob-
lems related to holomorphic functions. In particular,
such functions are naturally considered not only in
domains, but also on complex manifolds - smooth
manifolds of even real dimension, the neighbourhood
relations of which are biholomorphic. Among these,
Stein manifolds (cf. Stein manifold) - natural generali-
zations of domains of holomorphy - play a special
role
Several problems in analysis may be reduced to the
problem of constructing, in a given domain, a hol~
morphic function with given zeros or a meromorphic
function with given poles and principal parts of the
Laurent series. In the plane case, these problems have
been solved for arbitrary domains by the theorems of
Weierstrass and Mittag-Leffler and their generaliza-
tions. The spatial case is different - the solvability of
the corresponding problems, the so-called Cousin prob-
lems, depends on certain topological and analytic pro-
perties of the complex manifolds considered:
The key step in the solution of the Cousm problems
is to construct - starting from locally-defined func-
tions with given properties - a global function,
defined on the whole manifold under consideration and
having the same local properties. Such kinds of con-
structions are very conveniently effected using the
theory of sheaves, which arose from the algebrai~
topological treatment of the concept of an analytIc
function, and which has found important applications
in various branches of mathematics. The solution of the
Cousin problems by methods of the theory of sheaves
was realized by H. Cartan and J.-P. Serre.
B. V Shabat
The contemporary theory of analytic functions and their
generalizations. This is one of the most important
branches of analysis. it is closely connected with quite

176

diverse branches of mathematics and it has numerous
applications in theoretical physics, mechanics and tech-
nology.
Fundamental investigations on the theory of analytic
functions have been carried out by Soviet mathemati-
cians. Extensive interest in the theory of functions of a
complex variable emerged in the Soviet Union at the
begining of the 20-th century. This was in connection
with noteworthy investigations by Soviet scientists on
applications of the theory of analytic functions to vari-
ous problems in the mechanics of continuous media.
N.E. ZhukovskiI and S.A. Chaplygin solved very
important problems in hydrodynamics and aerodynam-
ics by using methods of the theory of analytic func-
tions. In the works of G.V. Kolosov and N.l. Muskhel-
ishvili these methods were applied to fundamental
problems in the theory of elasticity. In subsequent
years the theory of functions of a complex variable
underwent extensive development. The development of
various aspects of the theory of analytic functions was
determined by the fundamental research of, among oth-
ers, V.V. Golubev, N.N. Luzin, 1.1. Privalov and V.l.
Smirnov (boundary properties), M.A. Lavrent'ev
(geometric theory, quasi-conformal mappings and their
applications to gas dynamics), M.V. Keldysh, M.A.
Lavrent'ev and L.L Sedov (applications to problems in
the mechanics of continuous media), D.E. Men'shov
(theory of monogeneity), M.V. Keldysh, M.A.
Lavrent'ev and S.N. Mergelyan (approximation theory),
LN. Vekua (theory of generalized analytic functions
and their applications), A.O. Gel'fond (theory of inter-
polation), N.N. Bogolyubov and V.S. Vladimirov
(theory of analytic functions of several variables and its
application to quantum field theory). The development
of the theory of analytic functions of one and several
complex variables and their generalizations is continu-
ing. See Boundary properties of analytic functions;
Quasi-conformal mapping; Boundary value problems of
analytic function theory; Approximation of functions of
a complex variable.
References
[1] PRIWALOW. 1.1. [1.1. PRIVALOV]: Einfuhrung in die Funk-
tionentheorie, 1-3, Teubner, 1958-1959 (translated from the
Russian).
[2] SMIRNOV, V.l.: A course of higher mathematics, 111-2, Addison-
Wesley, 1964 (translated from the Russian).
[3] MARKuSHEVICH, A. I. : Theory offunctions of a complex variable,
1-3, Chelsea, 1977 (translated from the Russian).
[4] LAVRENT'EV, M.A. and SHABAT, B.V.: Methoden der komplexen
Funktionentheorie, Deutsch. Verlag Wissenschaft., 1967
(translated from the Russian).
[5] GoLUZIN, G.M.: Geometric theory offunctions of a complex
variable, Arner. Math. Soc., 1969 (translated from the Russian).
[6] EVGRAFov, M.A.: A nary tic functions, Saunders, Philadelphia,
1966 (translated from the Russian).
[7] SVESHNIKOV, A.G. and TnrnONOV, A.N.: The theory offunc-
tions of a complex variables, Moscow, 1967 (in Russian).
ANALYTIC FUNCTION
[8] FUKS, B.A.: Theory of anarytic functions of several complex vari-
ables, 1-2, Arner. Math. Soc., 1963-1965 (translated from the
Russian).
[9] VLADIMIROV, V.S.: Methods of the theory offunctions of several
complex variables, M.I.T, 1966 (translated from the Russian).
[lO] MARKUSHEVICH, A.I.: Notes on the history of the theory of ana-
rytic functiOns, Moscow-Leningrad, 1951 (in Russian).
[II] Mathematics in the USSR during 30 years, 1917-1947,
Moscow-Leningrad, 1948, pp. 319-414 (in Russian).
[12] Mathematics in the USSR during 40 years, 1917-1957, Vol. I,
Moscow, 1959, pp. 381-510 (in Russian).
[l3] SHABAT, B.V.: Introduction to complex anarysis, 1-2, Moscow,
1976 (in Russian).
[l4] BITSADZE, A.V.: Foundations of the theory of anarytic functions
of a complex variable, Moscow, 1972 (in Russian).
A.A. Gonchar
B. V Shabat
Editorial comments. Complex analysis, the theory of
analytic functions, has been an active field of research up to
the present time also in the West. Some of the history may
be found in the extensive article in the Encyclopaedia Bri-
tannica entitled Analysis, complex.
There is a large number of textbooks on complex analysis
besides those listed above. Some good current books are
listed below: [A1J - [A10J for the case of one variable, [A11J
- [A 18J for the case of several variables.
Some final comments on complex analysis of several vari-
ables are in order. The 'Osgood - Brown theorem on the
removability of compact singularities' is usually called the
Hartogs extension theorem in the West. The modern proof,
due to L. Ehrenpreis, makes use of the inhomogeneous
Cauchy - Riemann or a-equations. The so-called a-method,
extensively developed by J.J. Kohn and L. Hbrmander, has
become one of the three most powerful techniques available
in complex analysis today. Among other things, it can be
used to obtain solutions of the Levi problem (are domains of
holomorphy the same as pseudo-convex domains?) and the
Cousin problems, cf. [A 15]. These fundamental problems
had been settled previously by methods now belonging to
sheaf theory, and which go back to K. Oka (1936 - 1954),
and have culminated in the sheaf-cohomology theory of H.
Cartan, J.-P. Serre, H. Grauert and others, cf. the elegant
treatment in [A12J. The third and most recent technique
consists in the use of suitable integral representations as
developed by G.M. Henkin [G.M. KhenkinJ and E. Ramirez
de Arellano, cf. [A14].
References
[A1] AHLFORS, L.V.: Complex analysis, McGraw-Hill, 1966.
[A2] CONWAY, J.B.: Functions of one complex variable, Springer,
1973.
[A3] DINGHAS, A.: Vorlesungen uber Funktionentheorie, Springer,
1961.
[A4] HENRICI, P.: Applied and computational complex analysiS, 1-
3, Wiley, 1974-1986.
[AS] HILLE, E.: Analytic function theory, 1-2, Chelsea, reprint,
1974.
[A6] NEHARI, Z.: Conformal mapping, Dover, reprint, 1975.
[A7] NEVANILINNA, R.: Analytic functions, Springer, 1970
(translated from the German).
[A8] NEvANLINNA, R. and PAATERO, V.: Introduction to complex
analysis, Addison-Wesley, 1969 (translated from the Ger-
man).
[A9] RUDIN, W.: Real and complex analysis, McGraw-Hili, 1974.
177
ANAL YTIC FUNCTION

!
00
[A 10] lITCHMARSH, E.C.: The theory offunctions, Oxford Univ. fez) = CkZ- k /",
Press, 1939. k=m
[A11] GRAUERT, H. and FRITZSCHE, K.: Several complex variables,
Springer, 1976 (translated from the German). where v> I; the number v-I is called the branching
[A 12] GRAUERT, H. and REMMERT, R.: Theory of Stein spaces, order. The branched elements generalize the concept of
Springer, 1979 (translated from the German). an element of an analytic function, which in this con-
[A13] GUNNING, R. and ROSSI, H.: Analytic functions of several
complex variables, Prentice-Hall, 1965.
nection is also called an unramified (for v= 1) regular
[A14] HENKIN, G.M. [G.M. KHENKIN) and LEITERER, J.: Theory of (for m ~O) element.
functions on complex manifolds, Birkhauser, 1984. As the simplest element (D, /) of an analytic function
[A 15] HORMANDER, L.: An introduction to complex analysis in
fez) of several complex variables Z=(Zb ... ,Zn), n>l,
several variables, North-Holland, 1973.
[A 16) KRANTZ, S.G.: Function theory of several complex variables, one can take a multiple power series
Wiley (Interscience), 1982. OC!

[A17] NARASIMHAN, R.: Several complex variables, Univ. of Chi- fez) = ~ ck(z-a)k = (2)
cago Press, 1971. Ik 1=0
[A18] RUDIN, W.: Function theory in the unit ball ofC n , Springer,
1980.
[A19] RANGE, R.M.: Holomorphic functions and integral represen-
tations in several complex variables, Springer, 1986. where a=(aJ, . . , ,an) is the centre,
AMS 1980 Subject Classification: 30-XX, 32-XX I k I = k I + ... + k n , Ck =Ck, . . . Ck.,

(z -al =(z -ad' ... (z -anl', and D IS some

ANALYfIC FUNCfION, ELEMENT OF AN - The
collection (D, /) of domains D in the plane C of a com-
plex variable z and analytic functions f(z) given on D
by a certain analytic apparatus that allows one to effec-
tively realize the analytic continuation of f(z) to its
whole domain of existence as a complete analytic func-
tion. The simplest and most frequently used form of an
element of an analytic function is the circular element
in the form of a power series
polydisc
D = {ZEC": IZj-ajl<Rl'j=I, ... ,n}
in which the series (2) converges absolutely. However,
for n > lone has to bear in mind that a polydisc is not
the exact domain of absolute convergence of a power
series.
The concept of an element of an analytic function is
close to that of the germ of an analytic function.

fez) = ~ck(z-a)k (I) References

k =0
and its disc of convergence D={zEC: I z-a I <R}
with centre at a (the centre of the element) and radius of
convergence R >0. The analytic continuation here is
achieved by a (possibly repeated) re-expansion of the
series (1) for various centres b, I b-a I ~R, by formu-
las like OC!
fez) = ~dn(z-b)" = Co+
n =0
+[Cj(b -a)+Cj (z -b)]+
+[c2(b -a)2 +2C2(b -a)(z -b)+C2(Z -b)2]+ ....
Anyone of the elements (D,/) of a complete analytic
function determines it uniquely and can be represented
by means of circular elements with centres a ED. In the
case of the centre at infinity, a = 00, the circular ele-
ment takes the form
f(Z) = ~CkZ-k
k =0
with domain of convergence D = {z EC: I z I > R}.
In the process of the analytic continuation, f(z) may
turn out to be multiple-valued and there may appear
corresponding algebraic branch points (Gf. Algebraic
branch point), that is, branched elements of the form
-Yo
f(Z) = ~ck(z-a>""
k =m
[I] MARKUSHEVICH, A.I.: Theory offunctions of a complex variable,
1-2, Chelsea, 1977 (translated from the Russian).
[2] VLADIMIROV, V.S.: Methods of the theory offunctions of several
complex variables, M.I.T, 1966 (translated from the Russian).
E.D. Solomentsev
Editorial comments. For n> 1 the domain of absolute
convergence of a power series is a so-called Reinhardt
domain, cf. [A 1].
References
[AI] HORMANDER, L.: An introduction to complex analysis in several
variables, North-Holland, 1973, Chapt. 2.4.
AMS 1980 Subject Classification: 30A05, 30610,
30640, 32A30
ANALYTIC FUNCflONAL - An element f of the
space H' (Q), the dual of the space H(Q) of analytic
functions defined on an open subset Q of en, i.e. a
functional on H(Q). Thus, a distribution with compact
support is an analytic functional. There exists a com-
pact set K C Q, said to be the support of the analytic
functional f, on which f is concentrated: For any open
set w ~ K the functional f can be extended to H (w) so
that for all u EH (Q) the following inequality is valid:
If(u) I ~ Cwsup I u I,
where C w is a constant depending on w. There exists a
measure p, with support in K such that

178
 ANALYTIC GEOMETRY

f(u)=judp.. be the Cartesian orthogonal coordinates of the point M

OJ in the system Oxy. A point M with abscissa x and ordi-
An analytic functional is defined in a similar manner nate y is denoted by the symbol M(x,y). The coordi-
on a space of real-valued functions.
. -
nates of a Doint M clearlY determine its location with
M.l. Vottsekhovskit respect to the system Oxy.
Editorial comments. For applications to partial differen- Let a line L be given on a plane 'fT with a given
tial equations, see [A 1]. Cartesian orthogonal coordinate system Oxy. Using the
concept of coordinates of a point it is possible to intro-
References
[A 1] EHRENPREIS, L.: Fourier analysis in several complex vari- duce the concept of the equation of the line L with
ables, Wiley (Interscience), 1970. respect to system Oxy as an equation of the type
[A2] HORMANDER, L.: An introduction to complex analysis in F(x, y) = 0, which will be satisfied by the coordinates x
several variables, North-Holland, 1973, Sect. 4.5.
and y of any point M on L, but not by those of any
AMS 1980 Subject Classification: 46F15 point which does not lie on L.
The basic idea of the method of coordinates on a
ANALYTIC GEOMETRY - A branch of geometry. plane is that the geometric properties of the line L can
The fundamental concepts of analytic geometry are the be clarified by studying its equation F (x, y) = 0 by ana-
simplest geometric elements (points, straight lines, lytic and algebraic tools. Thus, for instance, the prob-
planes, second-order curves and surfaces). The principal lem of the number of intersection points between a
means of study in analytic geometry are the method of straight line and a circle is reduced to the analytic
coordinates and the methods of elementary algebra. problem of the number of solutions of the set of equa-
The genesis of the method of coordinates is closely tions of the line and of the circle.
linked with the intense development of astronomy, This was in fact the method of study of the proper-
mechanics and technology in the seventeenth century. ties of an ellipse, a hyperbola and a parabola, which are
In his Geometrie, R. Descartes (1637) gave a clear and the lines of intersection between a circular cone and
exhaustive account of this method and of the founda- planes which do not pass through its vertex (cf. Conic
tions of analytic geometry. P. Fermat, a contemporary sections).
of Descartes, was also familiar wih the principles of this The so-called algebraic curves of the first and second
method. Subsequent development of analytic geometry orders are a subject of systematic study in plane ana-
is due to the studies of G. Leibniz, I. Newton and, par- lytic geometry; in Cartesian orthogonal coordinates
ticularly, L. Euler. The tools of analytic geometry were these are described by algebraic equations of the first
utilized by J.L. Lagrange in his construction of analytic and second degree, respectively. The curves of the first
mechanics and by G. Monge in differential geometry. order are rectilinear and, conversely, any straight line is
In our own days, analytic geometry has no significance described by an algebraic equation Ax + By + C = 0 of
as an independent branch of science, but its methods the first degree. The curves of the second order are
are extensively employed in various fields of mathemat- described by equations of the form
ics, mechanics, physics and other sciences. AX2 + Bxy + Cy2 + Dx + Ey + F=O. The principal tech-
The principle of the method of coordinates is as fol- nique in the study and classification of such curves is
lows. Consider, for example, two mutually- to select a Cartesian orthogonal coordinate system in
perpendicular straight lines Ox and 0' in a plane 'fT. which the equation assumes its simplest form, after
These lines, including their direction, the coordinate which the simple equation is studied. Cf. Second-order
origin 0 and the selected scale unit e, form a so-called curve.
Cartesian orthogonal system Oxy of coordinates on the In the analytic geometry of space, the Cartesian
plane. The straight lines Ox and 0' are called the orthogonal coordinates are x, y and z (the abscissa, the
abscissa (axis) and the ordinate (axis) respectively. The ordinate and the applicate), and points M are described
location of any point M in the plane relative to this exactly as in plane analytic geometry. Any given sur-
system Oxy can be determined as follows. Let Mx and face S in space has its own equation F(x,y, z)=O in
My be the projections of M onto Ox and Oy, and let the coordinate system Oxyz, and its geometric proper-
the numbers x and y be the magnitudes of the segments ties are studied by studying its equation by analytic
OMx and OMy (the magnitude x of the segment OMx , and algebraic means. A curve L in space is given as the
for example, is equal to the length of this segment, intersection of two surfaces S \ and S 2. If the respective
taken with the plus sign if the direction from 0 to Mx equations of S \ and S 2 are F\ (x, y, z)=O and
is the same as the direction of the straight line Ox, and F 2 (x, y, z) = 0, these equations, considered together, are
taken with the minus sign in the opposite case). The the equation of the curve L. Thus, a straight line L in
numbers x (the abscissa) and y (the ordinate) are said to space may be considered as the intersection of two

179
ANALYTIC GEOMETRY
planes. In analytic geometry of space a systematic
study is made of the so-called algebraic surfaces of the
first and second orders. It was fQund that only planes
are algebraic surfaces of the first order. Surfaces of the
second order have equations of the type
Ax2 +By2 +Cz 2+ Dxy + Eyz +Fxz +
+Gx+Hy+Mz+N = O.
The principal method of study and classification of
such surfaces is to select a Cartesian coordinate system
in which the equation of the surface assumes its sim-
plest form, the latter being then studied. Cf. Surface of
the second order.
References
[1] DESCARTES, R.: La geometrie, Leiden, 1637.
[2] WIELEITNER, H.: Die Geschichte der Mathematik von Descartes
bis zum Hiilfte des 19. Jahrhunderts, de Gruyter, 1923.
[3] EFIMOV, N.V.: A short course of analytic geometry, Moscow,
1967 (in Russian).
[4] IL'IN, V.A. and POZNYAK, E.G.: Analytical geometry, Mir, 1984
(translated from the Russian).
[5] ALEKSANDROV, P.S.: Lectures on analytical geometry, Moscow,
1968 (in Russian).
[6] POSTNIKOV, M.M.: AnalytiC geometry, Moscow, 1973 (in Rus-
sian).
[7] BAKHVALOV, S.V., MODENOV, P.S. and PARKHOMENKO, A.S.:
Collection of problems on analytiC geometry, Moscow, 1964 (in
Russian).
E. G. Poznyak
Editorial comments.
References
[A 1] BOL, G.: Elemente der analytischen Geometrie, Vandenhoeck
& Ruprecht, 1948.
[A2] BORSUK, K.: Analytic geometry, PWN, 1969.
[A3] STRUIK, D.J.: Lectures on analytic and projective geometry,
Addison-Wesley, 1953.
[A4] TODD, I.A.: Projective and analytical geometry, Pitman and
Sons, 1947.
AMS 1980 Subject Classification: 51 NXX
ANALYTle GEOMETRY - The theory of analytic
spaces (cf. Analytic space). The term was introduced
by J.P. Serre [I] in analogy with algebraic geometry.
References
[I] SERRE, I.-P.: 'Geometrie algebrique et geometrie analytique',
Ann. Inst. Fourier (Grenoble) 6 (1955-1956), 1-42.
A.L. Onishchik
AMS 1980 Subject Classification: 32BXX, 32CXX
ANALYTIC GROUP - A set G which possesses at the
same time the structure of a topological group and that
of a finite-dimensional analytic manifold (over a field k
that is complete in some non-trivial norm, cf. Nonn on
a field) so that the mapping G X G~G defined by the
rule (x, y )~xy -I is analytic. An analytic group is
always Hausdorff; if k is locally compact, then G is
locally compact. If k is, respectively, the field of real,
complex or p-adic numbers, then G is called a real,
180
complex or p-adic analytic group, respectively. An
example of an analytic group is the general linear
group GL(n, k) of the vector space k n over k (cf.
Linear classical group) or, more generally, the group of
invertible elements of an arbitrary finite-dimensional
associative algebra with a unit over k. In general, the
group of k-rational points of an algebraic group,
defined over k, is an analytic group. A subgroup of an
analytic group G which is a submanifold in G is called
an analytic subgroup; such a subgroup must be closed
in G. For example, the orthogonal group
O(n, k)={gEGL(n, k): tgg=I} is an analytic sub-
group in GL(n, k). All closed subgroups of a real or p-
adic analytic group are analytic, and each continuous
homomorphism of such groups is analytic (Carlan's
theorems, [1]).
An analytic group is sometimes referred to as a Lie
group [1], but a Lie group is usually understood in the
narrower sense of a real analytic group [2], [3] (cf. Lie
group). Complex and p-adic analytic groups are called,
respectively, complex and p-adic Lie groups.
The Cartan theorems formulated above signify that
the category of real or p-adic analytic groups is a com-
plete subcategory in the category of locally compact
topological groups. The question of the extent to which
these categories differ, i.e. as to when a locally compact
group G is a real analytic or a p-adic analytic group,
can be exhaustively answered: If G is real analytic, it
must contain a neighbourhood of the unit without
non-trivial subgroups [5] - [9]; if it is p-adic, it must
contain a finitely generated open subgroup U which is
a pro-p-group and whose commutator subgroup is con-
tained in the set uP' of p2-th powers of elements in U
[10]. In particular, any topological group with a neigh-
bourhood of the unit that is homeomorphic to a
Euclidean space (a so-called locally Euclidean topological
group, [4]) is a real analytic group. In other words, if
continuous local coordinates exist in a topological
group, it follows that analytic local coordinates exist;
this result is the positive solution of Hilbert's fifih prob-
lem [5], [11].
If the characteristic of the field k is zero, the most
important method in the study of analytic groups is the
study of their Lie algebras (cf. Lie algebra of an
analytic group).
For infinite-dimensional analytic groups cf. Lie
group, Banach.
References
[I] SERRE, J.-P.: Lie algebras and Lie groups, Benjamin, 1965
(translated from the French).
[2] PONTRYAGIN, L.S.: Topological groups, Princeton Univ. Press,
1958 (translated from the Russian).
[3] CHEVALLEY, c.: Theory of Lie groups, I, Princeton Univ. Press,
1946 (translated from the French).
[4] HELGASON, S.: Differential geometry and symmetriC spaces,
Academic Press, 1962.

[5) Hilbert's problems, Moscow, 1969, pp. 101-115 (in Russian).
[6) GLEASON, A.M.: 'Groups without small subgroups', Ann. of
Math. (2) 56, no. 2 (1952), 193-212.
[7) MONTGOMERY, D. and ZIPPIN, L.: 'Small subgroups for finite
dimensional groups', Ann. of Math. (2) 56, no. 2 (1952), 213-
241.
[8) Y AMABE, H.: 'On the conjecture of Iwasawa and Gleason',
Ann. of Math. (2) 58, no. 1 (1953),48-54.
[9) YAMABE, H.: 'A generalization of a theorem of Gleason', Ann.
of Math. (2) 58, no. 2 (1953),351-365.
[10) LAZARD, M.: 'Groupes analytiques p-adiques', Publ. Math.
IHES 26 (1965).
[II) KAPLANSKY, I.: Lie algebras and locally compact groups, Chi-
cago Univ. Press, 1971.
VL. Popov
Editorial comments. In Western literature a connected
Lie group is often called an analytic group.
Cartan's theorems usually go back to J. von Neumann
(cf. [A1], [A2]).
References
[A1) NEUMANN, J. VON: Collected works, Vol. 1, Pergamon Press,
1961, pp. 134-148.
[A2) NEUMANN, J. VON: Collected works, Vol. 1, Pergamon Press,
1961, pp. 509-548.
AMS 1980 Subject Classification: 10K40, 22A05,
22E05, 57SXX, 58A05
ANALYTIC IMAGE - An extension of the concept of
a complete analytic function, obtained on considering
all possible elements of an analytic function in the form
of generalized power series (Puiseux series)
00 00
~a
~P
(z-z 0 ),/n '~P
~a z-,/n
(*)
J)=m
Here z is a complex variable, m is an integer and n is a
natural number. The series converge in the domains
1 z -Zo 1 <r and 1 z 1 >r>O, respectively. An analytic
image can be identified with the class of all elements of
the form (*) which are obtained from each other by
analytic continuation. The analytic image differs from
the complete analytic function by the addition of all
ramified elements of the form (*) with n > 1, which are
obtained by analytic continuation of its regular ele-
ments with n = 1. After the introduction of a suitable
topology, the analytic image is converted to the
Riemann surface of the given function.
References
[I) MARKUSHEVICH, A. I. : Theory of functions of a complex variable,
3, Chelsea, 1977, Chapt. 8 (translated from the Russian).
E.D. Solomentsev
AMS 1980 Subject Classification: 30840
ANALYTIC 'LANDSCHAFr', analytic relief, modulus
surface - The geometric image of the modulus f (z) 1 1
of an analytic function f(z), z=x+ry. The analytic
'landschaft' of a function f (z) is the surface over the
(x, y)-plane with z-coordinate 1f (z) I. An analytic
'landschaft' sometimes provides a good illustration of
the behaviour of specific functions. For the relief of the
ANALYTIC MANIFOLD
principal functions see [2].
References
[I) MARKUSHEVICH, A.I.: Theory offunctions of a complex variable,
1, Chelsea, 1977, Chapt. 2 (translated from the Russian).
[2) JAHNKE, E. and EMoE, F.: Tables offunctions with formulae and
curves, Dover, reprint, 1945 (translated from the German).
ED. Solomentsev
Editorial comments. Instead of the German word
'Landschaft', the English 'landscape' is also used.
AMS 1980 Subject Classification: 30AXX
ANALYTIC MANIFOLD - A manifold with an ana-
lytic atlas, The structure of an n-dimensional manifold
M over a complete non-discretely normed field k on a
topological space is defined by specifying an analytic
atlas over k on M, i.e, a collection of charts (cf. Chart)
with values in k n covering M, any two charts of which
are analytically related. Two atlases are said to define
the same structure if their union is an analytic atlas.
The sheaf ~ of germs of k-valued analytic functions is
defined on an analytic manifold. The class of ringed
spaces (M, ~) which results in this way is identical
with the class of smooth analytic spaces over k,
If k is the field of real numbers R one speaks of
real-analytic manifolds; if k is the field of complex
numbers C, of complex-analytic or simply complex man-
ifolds; if k is the field of p-adic numbers Qp' of p-adic
analytiC manifolds. Examples of analytic manifolds
include the n-dimensional Euclidean spaces k n , the n-
dimensional projective spaces over k, the affine and
projective algebraic varieties without singular points
over k, and Lie groups and their homogeneous spaces,
The concept of an analytic manifold goes back to B.
Riemann and F. Klein, but was precisely formulated
for the first time by H. Weyl [4] for the case of
Riemann surfaces, i.e. one-dimensional complex mani-
folds. At present (the 1970's) it is natural to regard
analytic manifolds as a special case of analytic spaces
(cf. Analytic space), which may be roughly described as
'varieties with singular points', The concept of an ana-
lytic space was introduced in the 1950's and became
the principal subject of the theory of analytic functions;
many fundamental results obtained for analytic mani-
folds could be successfully applied to the non-smooth
case. For an account of the general properties of ana-
lytic manifolds over an arbitrary field see [3].
There is a close relationship between the theories of
real-analytic and differentiable manifolds (cf, Differen-
tiable manifold), and also between the theories of real-
analytic and complex-analytic manifolds, Clearly, the
natural structure of a manifold of class Coo is defined
on each real-analytic manifold. It was shown by H.
Whitney in 1936 that the converse proposition is also
true: It is possible to define on any paracompact mani-
fold of class Coo an analytic structure over R which
181
ANAL YTIC MANIFOLD
induces the initial smooth structure. It follows from
Grauert's theorem on the imbeddability of a paracom-
pact analytic manifold over R in. a Euclidean space that
this analytic structure is unambiguously defined up to
an isomorphism (not necessarily the identity) [2].
A natural structure of a real-analytic manifold (of
double dimension) is defined on all complex manifolds
M. The answer to the converse problem - viz. whether
a complex structure on a given real-analytic manifold
exists and whether it is unique - has been given in the
simplest cases only. Thus, if M is a connected two-
dimensional real-analytic manifold, then a necessary
and sufficient condition for the existence of a complex
structure on M is paracompactness and orientability,
while the problem of classification of these structures is
identical with the classical moduli problem for
Riemann surfaces (cf. Moduli of a Riemann surface).
There is a classification of compact analytic surfaces
(i.e. of two-dimensional complex manifolds, cf. Analytic
surface), which gives a partial answer to the above
problem for four-dimensional real-analytic manifolds.
On the other hand it is possible, using topological
methods, to identify classes of real manifolds that do
not permit almost-complex or, a fortiori, complex struc-
tures. Such manifolds include the spheres S2k for
k=f:l, 3. A description of complex structures which are
sufficiently near to a given complex structure is given
by the theory of deformations of analytic structures (cf.
Defonnation), in which an important role is played by
Banach analytic manifolds - infinite-dimensional
analogues of analytic manifolds.
References
[11 BOURBAKI, N.: Elements of mathematics. Differentiable and
analytic manifolds, Addison-Wesley, 1966 (translated from the
French).
[21 NARASIMHAN, R.: Analysis on real and complex manifolds,
Springer, 1971.
[31 SERRE, J.-P.: Lie algebras and Lie groups, Benjamin, 1965
(translated from the French).
[41 WEYL, H.: Die Idee der Riemannschen Fiache, Teubner, 1955.
A.L. Onishchik
Editorial comments. A much related basic problem in
complex analysis is the question whether there are any
complex structures on projective space besides the usual
one (and inducing the same topology). For n = 1 this is
very classical (all Riemann surfaces of genus zero are iso-
morphic to Pb). For n =2, uniqueness of the complex
structure follows from combined work of F. Hirzebruch, K.
Kodaira [A4], and S.T. Yau [AS]. For n =3 one has that a
compact manifold that is bimeromorphically equivalent to a
Kahler manifold and that is also topologically q is analyti-
cally isomorphic to q [A6].
References
[A 11 WHITNEY, H.: Complex analytic varieties, Addison-Wesley,
1972.
[A2] BARTH, W., PETERS, C. and YEN. A. VAN DER: Compact com-
plex surfaces, Springer, 1984.
182
[A3] WELLS, R.O.: Differential analysis on complex manifolds,
Springer, 1980.
[A4] lIIRzEBRUCH, F. and KODAIRA, K.: 'On the complex projec-
tive spaces', J. Math. Pures Appl. 36 (1957), 201-216.
[A5] YAU, S.T.: 'Calabi's conjecture and some new results in
algebraic geometry', Proc. Nat. Acad. Sci. USA 74 (1977),
1789-1799.
[A61 I'ETERNELL, T.: 'A rigidity theorem for P3 (C)', Manuscripta
Math. 50 (1985),397-428.
AMS 1980 Subject Classification: 32C05, 58A07
ANALYTIC MAPPING, analytic morphism - A mor-
phism of analytic spaces considered as ringed spaces
(cf. Analytic space; Ringed space). An analytic map-
ping of a space (X, (!) x) into a space (X, (!) y) is a pair
(jO,fl), where
Io:X ~ Y
is a continuous mapping, while
II : 10 I ( (!) y) ~ (!) x
is a homomorphism of sheaves of rings on X. If the
spaces are complex, an analytic mapping is also called
a holomorphic mapping.
If (X, (!) x) and (Y, (!) y) are reduced analytic spaces,
the homomorphism f1 is completely determined by the
mapping fo and is the inverse mapping of the germs of
functions corresponding to fo. Thus, in this case an
analytic mapping is a mapping f: X ~ Y such that for
any x EX and for any CPE (!) f(x) one has cpof E (!) x.
A fibre of an analytic mapping
1= VoJd: (X, (!) x) ~ (Y, (!) y)
at a point y E Y is the analytic subspace
F'(y) = Vol (y), (!) x / II(my )(!) x 1!~I(y)
of the space (X, (!) x), where my E (!) y is the sheaf of
germs of functions that vanish at the point y. Putting
d(x) = dimxF'Vo(x)), XEX,
one obtains the inequality
dimx X ~ dimfo(x) Y +d(x). (*)
If X and Yare reduced complex spaces, then the set
Xl = {XEX: d(x)~l}
is analytic in X for any r~o.
An analytic mapping f=(jo,fd is called flat at a
point x EX if (!} X.x is a flat module over the ring
(!} Y,fo(x)' In such a case (*) becomes an equality. An
analytic mapping is called flat if it is flat at all points
x EX. A flat analytic mapping of complex spaces is
open. Conversely, if fo is open, Y is smooth and all
fibres are reduced, then f is a flat analytic mapping.
The set of points of a complex or a rigid analytic space
X at which an analytic mapping f is not flat is analytic
in X. If X and Yare reduced complex spaces, while X
has a countable base, then Y contains a dense
everywhere-open set over which f is a flat analytic map-

ping. If an analytic mapping
j:(X, ll1x)~(Y, ll1 y)
of complex spaces is flat, then the set of y E Y at which
the fibre r
1 (y) is not reduced or normal is analytic in
(X, ~ x).
Let f: X ~ Y be an analytic mapping of reduced
complex spaces. If dim X < 00, then there exists a
stratification
o = X( -1) C; X(O) C; ... C; X(r;) C; ... ,
where X(r) are analytic sets and X (r) = X for large r,
with the following property: Any point
xEX(r)\X(r-l) has a neighbourhood U in X such
that f(U nX(r is a local analytic set in Y, all irredu-
cible components of germs of which have dimension r
at f(x). If fis proper, then f(X) is an analytic set in
X. This is a particular case of the finiteness theorem for
analytic mappings.
Let X, Y be complex spaces and let X be compact.
Then it is possible to endow the set Mor (X, Y) of all
analytic mappings f: X ~ Y with the structure of a
complex space such that the mapping
Mor(X, Y)XX ~ Y,
which maps the pair (j, x) into f (x), is analytic. In
particular, the group of automorphisms of a compact
complex space X is a complex Lie group, acting analyt-
ically on X.
References
[I] REMMERT, R.: 'Projektionen analytischer Mengen', Math. Ann.
130 (1956), 410-441.
[2] REMMERT, R.: 'Holomorphe und meromorphe Abbildungen
komplexer Rliume', Math. Ann. 133 (1957), 328-370.
[3] STEIN, K.: Colloquium for topology, Strasbourg, 1954.
[4] FRISCH, J.: 'Points de plattitude d'une morphisme d'espaces
analytiques complexes', Invent. Math. 4 (1967), 118-138.
[5] FISHER, G.: Complex analytic geometry, Springer, 1976.
D.A. Ponomarev
AMS 1980 Subject Classification: 58A07, 32HXX
ANALYTIC MODEL OF A LANGUAGE, analyzing
model of a language - A type of mathematical construc-
tion used in mathematical linguistics to describe the
structure of natural languages. These constructions are
used in the formal modelling of basic linguistic
categories and of the process of linguistic research
itself; in other words, certain sets of 'disordered' data
about the language (or, more exactly, about speech) are
used to deduce certain information about the structural
mechanism of the language, i.e. about its grammar in
the wider sense of this word. The 'functioning' of such
a model does not always have the character of an effec-
tive construction, since the set of initial data need not
be a constructive object; in principle, this does not
render such models inferior.
ANALYTIC MODEL OF A LANGUAGE
In the most fully developed analytic models of
languages, the set of initial data is usually an object
which models a set of grammatical sentences of a
natural language, viz. some fonnal language in a given
alphabet (vocabulary) V. If L is a language in the voca-
bulary V and 'f/u, v E V [uxv EL=*uyv EL], then one
says that the string x can be substituted by the string y
with respect to V and L; if each one the strings x and y
can be substituted by the other with respect to V and
L, then one says that x and y can be mutually substi-
tuted with respect to V and L. The concept of mutual
substitutability has a simple linguistic meaning: If L is
understood to be the set of grammatical sentences of
some natural language, the mutually substitutable sets
are 'syntactically equivalent' (i.e. fulfilling the same
syntactic functions) word combinations. In particular,
if a single-symbol string a (in linguistic interpretation:
a word) is mutually substitutable with a string x of
length> 1, x is a 'potential constituent', i.e. may be a
constituent of the linguistically natural constituent
structures of grammatical sentences of this language (cf.
Syntactic structure); in such a case the string x is called
a first rank configuration of the language L with resultant
a. Thus, in the English language, the string 'uniformly
continuous' can be regarded as a first rank configura-
tion with resultant 'continuous'. However, first rank
configurations do not exhaust all 'potential consti-
tuents'; for example, the word combination 'continuous
function' is not a first rank configuration, since only
words such as 'function', 'derivative', etc., can be sub-
stituted by it, but it cannot itself be substituted' by any
of these words ('f (x) is a uniformly continuous func-
tion' is a correct statement, but 'f (x) is a uniform
function' is not). Therefore one introduces the following
definition: If r> 1 is a natural number and if the con-
cept of a rank i configuration of the language L has
been defined for each i = 1, ... , r -1, then a string x
of length > 1 is called a rank r configuration of the
language L with resultant a, where a E V, if the following
condition holds: If a can be substituted by x with
respect to V and L, if z 1XZ 2 ELand if z 1XZ 2 does not
contain occurances of configurations of ranks lower
than r overlapping the considered x but not entirely
contained in it, then z I az 2 EL. In English, the word
combination 'continuous function' can be regarded as a
second rank configuration, the word 'function' being
the resultant (the word 'derivative', for example, is also
a resultant). It may be shown that, in a certain sense, a
language is fully determined by the set of its configura-
tions.
The configurational model belongs to the so-called
syntagmatic analytic models of a language, which are
intended to describe the relations between the elements
of segments of speech (in linguistics such relations are
183
ANALYTIC MODEL OF A LANGUAGE
called syntagmatic). Another class of analytic models
of a language are paradigmatic models, intended to
describe paradigmatic relations, i.e. relations between
the elements of a language in its system. In such
models the usual procedure is to construct certain rela-
tions on the vocabulary, which are often (but not
always) equivalence relations. In particular, paradig-
matic analytic models of a language serve to construct
formal analogues of traditional linguistic categories
such as parts of speech, case, gender, phoneme, etc. A
set of grammatical sentences also serves as the 'starting
material' in many paradigmatic analytic models of a
language; the concept of substitutability is employed in
several models of this type. The simplest equivalence
relation obtained in this way is the mutual substituta-
bility relation on the vocabulary (i.e. on the set of all
one-element strings); the equivalence classes induced by
this relation are called families. If one introduces an
additional relation on the vocabulary, consisting of 'to
be different forms of the same word' (more precisely, of
one lexeme; this is, e.g., the relation between 'to limit'
and 'limiting' or between 'number' and 'numbers'; it
involves a certain idealization in that it is assumed that
a word may be the form of one lexeme only; the
respective equivalence classes are called neighbour-
hoods), then these two relations may be used to intro-
duce certain other classifications, which may be
regarded as approximations to the formal analogues of
parts of speech and other traditional grammatical con-
cepts, in particular the concepts of case and gender of
nouns. These last two concepts are particularly inten-
sively studied; a number of models based on the con-
cept of grammaticality as well as other types of models
have been suggested for their formalization. In one
model, in particular, each case of nouns is treated as a
set of 'uniformly directed' forms of nouns, while each
gender is treated as a set of 'uniformly directing' nouns.
The initial data in this model, which serve to determine
the gender, are the vocabulary, the set of neighbour-
hoods, the set of nouns and the binary relation of
'potential subordination' in the vocabulary: a poten-
tially subordinates b if, in some 'linguistically natural'
dependency structure (cf. Syntactic structure) of a
grammatical sentence in this language, some occurrence
of b (immediately) depends on some occurrence of a;
additional initial data of the same nature are added to
determine the gender. Models based on 'potential
subordination' or other concepts related to dependency
structures apparently permit a better formalization of
traditional grammatical categories than do models
involving strings only, since in the former .models syn-
tactic and linear relations are separated (cf. Grammar,
transformational).
The mathematical apparatus employed in the con-
184
struction of analytic models of a language is usually
fairly simple. In some models it merely involves the
simplest concepts of set theory (such models are fre-
quently called set-theoretic, and this term is sometimes
extended to cover all analytic models of a language). In
other cases algebraic concepts are employed, in particu-
lar those of the theory of semi-groups and of the alge-
bra of binary relations; this is why the theory of ana-
lytic models of a language is sometimes referred to as
algebraiC linguistics.
References
[1] KULAGINA, O.S.: 'On a means of defining grammatical con-
cepts on the base of set theory', Problemy Kibemet. 1 (1958),
203-214 (in Russian).
[2] USPENSKIi, V.A.: 'On the Kolmogorov definition of case',
Byull. Obshch. Problemy Mashinnogo Perevoria, no. 5 (1957),
11-18 (in Russian).
[3] USPENSKIi, V.A.: 'A model for the concept of phoneme',
Voprosy Yazykoznaniya, no. 6 (1964), 39-53 (in Russian).
[4] MARKus, S.: Algebraic linguistics: analytical models, Acad.
Press, 1967.
[5] REvzIN, I.I.: Modelling and typologizing Slavic languages, Mos-
cow, 1967 (in Russian).
[6] GLADICU, A.V. [A.V. GLADIaij and MEL'i':UK, LA. [IA
MEL'CHUK]: Elements of mathematical linguistics, Mouton
Pub!., 1983.
[7] GLADKIi, A.V.: Formal grammars and languages, Moscow, 1973
(in Russian).
A. V. Gladkit
AMS 1980 Subject Classification: 68FXX
ANALYTIC NUMBER THEORY - A branch of
number theory. Analytic number theory deals with the
problems of distribution of primes, studies the
behaviour of number-theoretic functions, and the
theory of algebraic and transcendental numbers.
Distribution of prime numbers. a) The problem of the
distribution of prime numbers (primes) is one of the
most interesting and most difficult problems in analytic
number theory. The first result on the problem of the
distribution of primes was obtained by Euclid: There
are infinitely many primes. Let ?T(x) be the number of
primes not exceeding x. Euclid's theorem can then be
formulated as follows: ?T(x)~+ 00 as x~oo. The next
step was due to P.L. Chebyshev (1850), who proved
that:
1) The quantity ?T(x) satisfies the inequalities
ax bx
-I-
nx
< 7T(X) < -Inx
-,
where a~(ln2)j2; b~2In2.
2) If there exists a limit of ?T(x)ln(x)j x as x~+oo,
this limit is one.
The problem of the existence of this limit was solved
in 1896 by 1. Hadamard and by Ch.1. de la Vallee-
Poussin, who established thereby that
X
7T(X) -(I + 0(1)).
= -Inx

De la Vallee-Poussin demonstrated a much more gen-
eral assertion, viz. if
x d
7T(X) = [ Inll + R(x),
then
R(x) = O(xe-aVm-;),
where a>O is an absolute constant (cf. de la Vallee-
Poussin theorem). This problem was solved by methods
of the theory of functions of a complex variable. The
problem of estimating R (x) is closely connected with
the problem on the behaviour of a certain function of a
complex variable, which was first (1859) studied by B.
Riemann and which is now known as Riemann's zeta-
function. This function is defined by
res) = 1- 5+2- S + .,. +n- S + "', s=o+it, 0>1.
Even earlier (1737, 1749) the zeta-function with real s-
values was studied by L. Euler, who demonstrated the
identity specifying the connection between r(s) and
prime numbers:
res) = (1-2- )-I(1-3- )-1
S S '" (l_p-s)-I " ' , 0>1,
where the product is taken over all prime numbers. The
function t(s), which is given in the form of a series for
0'> 1, can be analytically extended to the entire plane
of the complex variable; the resulting function is
analytic in the entire complex plane except at the point
s= 1, where it has a simple pole with residue one. The
problem of estimating the remainder R (x) in the
asymptotic formula for the distribution of primes is
closely connected with the problem of the distribution
of the zeros of t(s) in the 'critical' strip 0..;;0''';; 1.
According to the Riemann hypothesis, all the zeros of
t(s) in the critical strip lie on the straight line 0' = 1 / 2.
It follows from this hypothesis that R (x) = O( Vx In x).
Conversely, the Riemann hypothesis on the zeros of
t(s) follows from the relation R (x) = O(x 1/2+<), where
>0 may be as small as one pleases (cf. Riemann
hypotheses). Hadamard and de la Vallee-Poussin
obtained the asymptotic law for the distribution of
primes by proving that t(s) cannot have any zeros if
O'~ 1. So-called ~-theorems were proved for R (x):
There exist two sequences X ~ + 00, y ~ + 00 such that
R(X) < _XI/2 InlnlnX,
hllnX
R(Y) > yl/2 I~ln~n/.
b) A second problem on the distribution of primes is
to estimate the difference between two successive
primes, i.e. the number I:1n = Pn + 1 - Pn, where pn is the
n-th prime. Here too, the first general result is due to
Chebyshev, who showed that at least one prime number
is contained between Nand 2N, where N~ 1 (the Ber-
trand postulate). The estimate of I:1n is closely connected
ANALYTIC NUMBER THEORY
with the function N(a, T), which is the number of zeros
of t(s) in the rectangle 1/2";;a..;;0'..;;1, It I..;;T. The
function N(a, T) is in turn closely connected with the
function r(1 /2 + it). The existing hypotheses are:
N(a, T)=0(T 2(1-a)+<) (the density hypothesis) and
r(1 /2 + it)= O(t<) (the Lindelof hypothesis), where >0
may be as small as one pleases. Lindelofs hypothesis
follows from the Riemann hypothesis on t(s); the den-
sity hypothesis follows from Lindelofs hypothesis, and
it follows from the density hypothesis that
I:1n = O(n 1/2+<). It was proved that I:1n = O(nY), where
y<1.
c) The problem of the distribution of primes in arith-
metical progressions nk + I, (k, 1)= 1, n =0, 1, ... ,
leads to the problem on the zeros of special zeta-
functions - the so-called Dirichlet L-series, of the
form
all-s+a22-5+ ... +ann- s + ...
for 0'> 1, where an depends on n and on the difference
k of the series (Dirichlet characters modulo k).
The problem on the distribution of zeros of Dirichlet
L-series and the distribution of primes in arithmetical
progressions have their own specific features. One of
the principal achievements in this field, due to c.L.
Siegel (1935), is as follows: Let 7T(X, k, I) be the number
of primes not exceeding x in the series nk + I,
n =0, 1, .... Then
7T(X, k, /) -- I (dt
CP(kH TrV + O(x In -A x,)
where <P(k) is the Euler function and l..;;k..;;In B x,
where A and B are arbitrary fixed numbers. Informa-
tion about the distribution of prime numbers in arith-
metical progressions is extensively used in solving addi-
tive problems involving prime numbers. See also
Distribution of prime numbers.
Additive problems. Additive problems in analytic
number theory include problems involving a special
type of integer equations. The main problems of this
kind are: To prove the solvability of a given equation
and to find an asymptotic formula for the number of
solutions of a given equation. The latter problem is
much more difficult and, in a certain sense, its positive
solution constitutes an answer to the first problem.
Classical examples of additive problems are the Waring
problem, the Goldbach problem and the
Hardy - Littlewood problem.
Waring's problem (1770) can be formulated as fol-
lows: Let J k n(N) be the number of solutions in positive
integers XI, . , Xn of the equation
x7 + ... +xZ = N, (1)
where N ~ 1 is an integer. It is required to prove that
there exists a number ko =ko(n), where ko depends
185
ANAL YTIC NUMBER THEORY
only on n, such that Jk,n(N)~l if k~ko. In other
words, it is required to prove that any number N ~ 1
can be represented by a sum of n-th powers of positive
integers, the number of terms in this representation
depending only on n. J.L. Lagrange (1770) solved this
problem for n = 2 by showing that any positive integer
is the sum of four squares of integers. D. Hilbert (1909)
was the first to give a general solution to Waring's
problem. G.H. Hardy and J.E. Littlewood (1924) used
their circle method to prove that for hn(N), if
k~n2n + 1, an asymptotic formula of the form
Jk,n(N) = ANk/n-l +O(Nk/n-1-y) (2)
is true, where y>O and A =A (N)~c>O; c is an abso-
lute constant. And since there is an infinitely large
number of numbers N which, for k = n, are not the sum
of n-th powers, i.e. h,n(N)=O, there arose the problem
of establishing the true order of k, as a function of n,
for which equation (1) is solvable and formula (2) is
valid. The strongest results on this problem were
obtained by I.M. Vinogradov (1934) who showed that:
a) h,n(N)~1 for N~No(n) if k~3n Inn + lln;
b) formula (2) is true for k~4n2lnn.
A second classical additive problem - the
Goldbach - Euler problem (1742) - is as follows: Let
J(N) be the number of solutions in primesPI,P2,P3 of
the equation P 1 +P2 +P3 =N; prove that for odd N~7,
J (N)~ 1. It was shown in 1937 by Vinogradov that (an
asymptotic formula for J(N:
J(N) = BN 2 1n- 3 N +O(N 2 In- 35h N),
where B = B(N)~0.6. It follows, in particular, that
J (N) ~ 1 for N ~ No, i.e. a solution of the
Goldbach - Euler problem for sufficiently large N.
Additive problems also include the
Hardy- Littlewood problem (1923): Each N~2 can be
represented in the form N=p+x 2 +y2, where p is a
prime number and x and yare positive integers. Yu.V.
Linnik proved in 1958 that if WeN) is the number of
solutions of this equation, the asymptotic formula
I:: +
WeN) = O(N In -1028 N),
where a=a(N)~cI >0 and CI is an absolute constant,
is true, It follows that W(N)~ 1 if N~N 0, which is the
solution of the Hardy- Littlewood problem for suffi-
ciently large N. There are many additive problems
which have remained unsolved for hundreds or even
thousands of years, These include, for example, the
problem on the infinite number of twin primes, i,e.
pairs of primes P and q such that I p - q I = 2, the
binary problem of Goldbach - Euler, viz. that any even
number ~4 is a sum of two primes; and "the problem
on the existence of an infinite number of primes in the
sequence of the form n 2 + 1. See also Additive problems,
186
Behaviour of number-theor~tic functions. The theory of
numbers comprises several classical functions: .p(n) -
the number of numbers not exceeding n and mutually
prime with n (the Euler function); T(n) - the number
of divisors of n; /LCn) - the Mobius function; A(n) -
the Mangoldt function, etc. Even though all these func-
tions behave fairly 'irregularly', their average values can
properly be studied. The average value of a function
f (n) is understood to mean the value of L.n,,;;,J
(n) / x.
The problem of estimating the average value of J.L(n) is
equivalent to the problem of bounds on the zeros of the
Riemann zeta-function. The problem of the asymptotic
behaviour of the average value of A(n) is equivalent to
the problem of the asymptotic formula for '/T(x), i.e.
again to the problem of bounds on the zeros of the
Riemann zeta-function. The results obtained on all
these problems are the same as those obtained in the
problem of the distribution of prime numbers. This
does not apply to the problem of the asymptotic
behaviour of the average value of 'T(n) or, to put it in a
somewhat different manner, the problem of the asymp-
totic formula for the sum of the values of 'T(n). Let
<II(N) = 7'(1)+ .. , +7'{N),
P(N) is then the number of integer points below the
hyperbola y = N / x. Thus, to find the asymptotic
behaviour of P(N) amounts to finding the asymptotic
behaviour of the number of integer points in expanding
domains. This class of problems also includes the prob-
lem of the integer points in a disc, i.e. the problem of
determining the number
G(N) = 2 ~ 1,
x +y .;;;N
where x and yare integers, and of extending these
problems to arbitrary domains, both in a plane and in
space. It was shown by P. Dirichlet in 1849 that
<II(N) = N In N +(2y+ I)N + R 1(N),
where R 1(N) = OeVN); it was also shown by C.F.
Gauss in 1863 that
G(N) = 'TTN + RiN),
where R 2 (N) = ocvii). The problems of finding the
best possible estimates for R I (N) and R2(N) are
known as divisor problems and the circle problem.
respectively, G.F. VoronoI (\ 903) obtained the formula
Rl(N) = O(NI3InN),
while the formula found by W. Sierpiilski (1906) was
R 2 (N) = O(N1 1).
In addition, the Q-theorems:
RI(N) = [2(NI/4) and R 2 (N) = [2(NI/4)
have also been demonstrated. The estimates of R I (N)

and R2(N) obtained in our own days (1976) are some-
what better than those of Vorono! and Sierpiilski.
A problem related to those just discussed is the
asymptotic behaviour of the sum of the fractional parts
of various types of functions or an equivalent problem,
viz. the problem of the distribution of the fractional
parts of various types of functions. Let {a} be the frac-
tional part of the number a. Then, if F(x) is a real
function, the problem concerns the asymptotic
behaviour of the two functions
T](N) = ~ {F(n)},
n<O;,N
~ I.
If for any O<y';;;; I
T 2 (y; N) = "IN +o(N),
one says that the fractional parts of F(x) are uniformly
distributed. The uniform distribution of the fractional
parts of F(n) can also be expressed in terms of the
asymptotic behaviour of T I (N). The first results on the
uniform distribution of the fractional parts of polyno-
mials, and criteria for uniform distribution, were found
by H. Weyl in 1916. The most accurate results in this
field were obtained by Vinogradov, who also found
asymptotic formulas for Tl (N) and T2(N) in cases in
which n runs through a part of the set of integers not
larger than N, in particular through the set of prime
numbers. Not much is known about the distribution of
the fractional parts of functions which grow faster than
polynomials. For instance, nothing is known about the
distribution of the fractional parts of the function
(3/2Y
Algebraic and transcendental numbers. The theory of
algebraic and transcendental numbers includes prob-
lems connected with the arithmetical nature of given
numbers or classes of numbers. Consider polynomials
with integer coefficients and with leading coefficient
one; if a is a root of such a polynomial of degree nand
is not a root of a polynomial of a lower degree, it is
said to be an algebraic number of degree n; if n = 1, the
number a is said to be rational. If, on the other hand, 0:
is not an algebraic number, it is said to be transcenden-
tal. There are 'much fewer' algebraic than transcenden-
tal numbers; 'almost any' number is transcendental,
but determining whether a specific number is transcen-
dental or algebraic is very difficult. The principal
'characteristic' of algebraic numbers is the fact that
they are very 'poorly' approximated to by rational
numbers. This assertion (Liouville's theorem, cf.
Liouville theorems, 1844) is formulated as follows: If 0:
is an algebraic number of degree n, then
Ia-~ I >
cq-.,
ANAL YTiC NUMBER THEORY
where ,,=n, c =c(o:O is a constant which depends
only on a, while p and q are arbitrary integers. The
next decisive step in this problem was made by A.
Thue (1909) whose ideas exerted a major influence on
the theory of transcendental numbers. He showed that
,,=n /2+, >0. Subsequently, the value of " was
reduced by several workers, and it was shown by K.F.
Roth in 1955 that ,,=2+ (it is known that ,,;;;;'2). A
drawback of all these theorems (except for Liouville's
theorem) consists in the fact that they are ineffective,
i.e. c .;annot be computed from known 0: and .
Approximation problems are connected with a cer-
tain class of problems in the theory of indefinite equa-
tions. E.g., Thue employed his approximation theorem
to prove that the number of integral solutions of the
equation Fn(x,y)=a, where Fn(x,y) is a form with
integral coefficients of degree n;;;;.3 and a is a non-zero
integer, is finite (this theorem is also ineffective, i.e. no
bound on the number of solutions of the equation can
be given).
Another approach employed in this theory are tran-
scendence proofs of numbers. The first results on this
subject were obtained towards the end of the
nineteenth century by Ch. Hermite (1873), who proved
that the number e is transcendental, and by F. Lin-
demann (1882), who proved that the number 7T is tran-
scendental and thus gave a negative solution to the
problem of squaring the circle. According to the
theorem of Gel'fond - Schneider (1934), a P is a tran-
scendental number if 0: is an algebraic number, 0:=f=0, 1
and f3 is an algebraic number of degree ;;;;.2 (Hilbert's
seventh problem). Beginning in 1967, A. Baker
obtained a number of effective theorems on the estima-
tion of linear forms in logarithms of algebraic numbers.
As a result, Thue's theorem on the number of represen-
tations of an integer by a form was effectively demon-
strated. There are still many unsolved problems in the
theory of transcendental numbers; these include the
transcendental nature of Euler's constant
"I = lim
n--+oo
[Inn -1 _1--
2
... _1-],
n
the problem of algebraic independence of the numbers
e and 7T, etc.
Some methods used in analytic number theory. a) The
method of complex integration. This method grew out of
Euler's method of generating functions which is often
used for the solution of problems in elementary
mathematics. The method is based on the following for-
mula (a discontinuous multiplier):
if x>l,
_1 f E.ds = if X = 1,
27Ti Res =u>O S 2
o ifO<x<l,
187
ANAL YTIC NUMBER THEORY
where the integral is taken over the straight line
Res =a>O, s =a+it. E.g., for Res> 1,
f(s) = -tf};) = n~IA(n)n-s,
and for x =N + 1 /2 one obtains
~ A(n) = _1_. f tiSil ds.
n';;;x 2Wl Res=2 S
The left-hand side of this equation is the Chebyshev
function, finding the asymptotic behaviour of which is
equivalent to the problem of prime numbers. The
right-hand side, after the separation of the main factor,
will be the smaller, the farther to the left the integra-
tion contour can be shifted. (Cf. Complex integration,
method of.)
b) The circle method (of
Hardy - Litlewood - Ramanujan). This method is used
mostly in additive problems. Below the application and
the nature of the circle method is considered in the
example of applying Vinogradov's trigonometric sums
to the ternary problem of Goldbach - Euler. Let m be
an integer. One then has (a discontinuous multiplier):
I
f e27r;am da = {I if m=O,
if m=f=O.
o
Therefore
I (1937) was the solution of several additive problems on
leN) = fS3(a)e2'1TtaN da, Sea) = ~ e'1T;XP,
o p,;;;N
and I (N) is the number of solutions in primes of the
equation PI +P2 +P 3 = N. The integration interval is
then subdivided into two parts - the principal and the
complementary interval; the principal interval contains
all the intervals of the form
!a-E-!
q
< _I,
q'
where (a, q)= 1, O",;;;;a",;;;;q, q",;;;;ln lO N, T=N In- 2o N; all
other intervals are within the complementary interval.
The principal intervals do not intersect. Moreover, for
a in the principal interval, the sum Sea) is 'close' to the
rational sum S(a / q). However, for 'small' q, q",;;;;lnN,
the distribution law of the prime numbers in a series
with difference q is known (e.g. Siegel's theorem), i.e.
the asymptotic behaviour of the sums S(a / q) is
known. In this way the principal term of the problem is
separated, which is the fundamental idea of the circle
method. If one now makes a non-trivial estimate of
I S(a) I on the complementary intervals (d. Vinogradov
method), the asymptotic formula in
Goldbach - Euler problem is obtained. (See also Circle
method.)
c) The method oj trigonometric sums. Most problems
in analytic number theory can be expressed by tri-
gonometric sums, i.e. finite sums of the form
188
S = ~ e 27r;F(x" ... x,), (3)
XIo'" ,x,.
where
x], . .
is a real function, while
F(x], ... , x n )
run through a set of P integers. Thus, the
,Xn
essence of many problems is the study of such sums, in
particular the problem of finding the most acurate pos-
sible estimate of the modulus of such a sum. A trivial
estimate of the sum (3) is P. The problem is to obtain
an estimate of the form
I S I ".; ; tu>,
where O",;;;;~",;;;; 1 is called the reducing factor. The first
non-trivial estimates of trigonometric sums for the case
when F=F(x) is a polynomial, while x=l, ... ,P,
were obtained in 1919 by Weyl, who at the same time
deduced a criterion for the uniform distribution of the
fractional parts of functions in terms of trigonometric
sums. The method of trigonometric sums was created
by Vinogradov who, by utilizing intrinsic arithmetical
properties of such sums, obtained very strong estimates
of the modulus for a broad class of these sums. He was
able to obtain in this way several fundamental results,
close to the limit of the possible, in several problems of
number theory (the Waring
Hilbert- Kamke problem, and Weyl sums, cf. Weyl
problem, the
sum). Another consequence of Vinogradov's method
prime numbers, in particular the Euler - Goldbach
problem. The fundamental idea in Vinogradov's
method was 'smoothing' (exponentiation of the tri-
gonometric sum and reduction of the estimate to the
theorem of averages for the estimation of Weyl sums,
the introduction of double trigonometric sums for the
estimation of sums of primes, etc.). (See also Tri-
gonometric sums, method of.)
d) The dispersion method and the method oj the large
sieve. Linnik in 1958 - 1960 proposed a dispersion
method for the solution of several additive problems in
number theory. He solved the Hardy- Littlewood
problem, the Titchmarsh problem on divisors and the
additive divisor problem. The basic concept in the
method is the dispersion of the number of solutions of
a given auxiliary equation related to the principal equa-
tion (see also Dispersion method). Fundamental results
were recently obtained by the method of the large sieve
proposed by Linnik in 1940 for the solution of the
problem of the smallest quadratic residue.
e) Methods in algebraic and transcendental number
theory. In order to prove his theorem on the approxi-
the mation of an algebraic number by a rational fraction,
Thue (cf. Thue method) constructed the polynomial
/(x.y) = (y -a)/l (x.y)+(x -a)"h(x, a)
with integer coefficients, where JI and 12 are also poly-
nomials. On the assumption that a is 'well' approxi-

mated by fractions PI / ql and P2 / q2 with sufficiently
large ql and q2, putting m~lnql /lnq2, and proving
that f(x, y) does not vanish if x =p I / ql, Y =P2 / q2, a
contradiction is obtained.
A.O. Gel'fond's proof of the transcendental nature of
numbers involves the function
Ina
f(z) = qo ~
~ 2,. ck./e(>.k+I)z, A= --.
k=OI=O Inf3
On the assumption that a is an algebraic number and
using Dirichlet's box principle, non-zero integers Ck,l
are chosen so that f (z) and 'many' of its derivatives
have 'many' zeros. A 'large' number of zeros will yield
'good' estimates from above for a 'large' number of
derivatives and points, whence, using the estimates
from below obtained by Liouville's theorem, it follows
that f (z) and 'many' of its derivatives have more zeros
than before. A repetition of this process results either
in A being a rational number or in zero values for ck,l,
which is in contradiction with their choice. See also
Algebraic number; Transcendental number.
References
[1] VINOGRADOV, I.M.: Selected works, Springer, 1985 (translated
from the Russian).
[2] VINOGRADOV, I.M.: Elements of number theory, Dover, reprint,
1954 (translated from the Russian).
[3] VINOGRADOV, I.M.: The method of trigonometric sums in the
theory of numbers, Interscience, 1954 (translated from the Rus-
sian).
[4] GEL'FOND, A.O.: Transcendental and algebraic numbers, Dover,
reprint, 1960 (translated from the Russian).
[5] DELONE, B.N.: The Peterburger school in number theory,
Moscow-Leningrad, 1947 (in Russian).
[6] KARATSUBA, A.A.: Fundamentals of analytic number theory,
Moscow, 1975 (in Russian).
[7] LINNIK, YU.V.: The dispersion method in binary additive prob-
lems, Amer. Math. Soc., 1963 (translated from the Russian).
[8] CHUDAKOV, N.G.: Introductions to the theory of Dirichlet L-
functions, Moscow-Leningrad, 1947 (in Russian).
[9] PRACHAR, K.: Primzahlverteilung, Springer, 1957.
[10] DAVENPORT, H.: Multiplicative number theory, Springer, 1980.
[II] TnCHMARSH, E.C.: The theory of the Riemann zetajunction,
Clarendon, 1951.
[12] HUA, L.-K.: 'Abschatzungen von Exponentialsummen und ihre
Anwendung in der Zahlentheorie', in Enzyklopaedie der
Mathematischen WissenschaJten mit Einschluss ihrer Anwen-
dungen, Vol. I, 1959. Heft 13, Teil I.
[13] BAKER, A.: 'Linear forms in the logarithms of algebraic
numbers (II)', Mathematika 14 (1967), 102-107.
[14] VINOGRADOV, A.I.: 'The density hypothesis for Dirichlet L-
series', Izv. Akad. Nauk SSSR Ser. Mat. 29, no. 4 (1965),903-
924 (in Russian).
[15] BOMBIERI, E.: 'On the large sieve', Mathematika 12 (1965),
201-225.
A.A. Karatsuba
Editorial comments. Liouville's theorem is closely
related with the theory of continued fractions (cf. Continued
fraction).
For a description of various results obtained on additive
problems and the distribution of prime numbers in arithmetic
progressions see Large sieve; Distribution of prime
numbers; Sieve method.
ANALYTIC POLYHEDRON
References
[A1] VAUGHAN, R.c.: The Hardy-Littlewood method, Cambridge
Univ. Press, 1981.
[A2] HALBERSTAMM, H. and RICHERT, H.E.: Sieve methods, Acad.
Press, 1974.
[A3] IVIt, A.: The Riemann zeta-function, Wiley, 1985.
AMS 1980 Subject Classification: 10-XX
ANALYTIC OPERATOR at a point Xo - An operator
A, acting from one Banach space into another, that
admits a representation of the form
00
A(xo +h)-Axo = ~ Ckhk,
k=l
where Ck is a form of degree k and the series converges
uniformly in some ball II h II <r. An operator is called
analytic in a domain G if it is an analytic operator at all
points of this domain. An analytic operator is infinitely
differentiable. In the case of complex spaces, analyti-
city of an operator in a domain is a consequence of its
differentiability (according to Gateaux) at each point of
this domain. Examples of analytic operators are
Lyapunov's integro-power series, and the Hammerstein
and Urysohn operators with smooth kernels on the
space C of continuous functions.
References
[I] HILLE, E. and PHILLIPS, R.: Functional analysis and semi-
groups, Arner. Math. Soc., 1957.
[2] KRASNOSEL'sKIi, M.A., ET. AL: Approximate solution of operator
equations, Woiters-Noordhoff, 1972 (translated from the Rus-
sian).
P. P. Zabrelko
AMS 1980 Subject Classification: 46G20, 58C10,
58C20, 47HXX
ANALYTIC PLANE, complex-analytic plane - A non-
empty set of points z = (z I, . . . , zn) in the complex vec-
tor space en that satisfies a system of equations
n
~aijzi = bJ , j= I, ... ,k; aij, bjEC;
i=1
rank II aij II = k < n.
The number k is called the complex codimension, while
n - k is called the complex dimension of the analytic
plane. The real dimension of the analytic plane equals
2(n -k) and is even, but not all even-dimensional real
planes in R2n = en are analytic planes. Complex one-
dimensional analytic planes are sometimes called com-
plex, or analytic, straight lines. See also Analytic surface.
E.D. Solomentsev
E.M. Chirka
AMS 1980 Subject Classification: 15A06, 51 N1 0,
32AXX
ANALYflC POLYHEDRON - A domain II of the
complex space en, n ~ 1, which can be represented by
inequalities 1J;(z) 1<1, where the functions J;(z),
189
ANALYTIC POLYHEDRON
i = I, ... ,m, are holomorphic in some domain D c en
containing II, i.e. II={zED: 1J;(z) 1<1,
i = I, ... ,m}. It is also assumed that II is compact in
D. If J;(z) are polynomials, the analytic polyhedron is
said to be a polynomial polyhedron. If m = nand
J;(z)=a;z;, the analytic polyhedron is called a polydisc.
The sets O';={zED: 1J;(z) I =1; 1Jj(z) I <l,j:t=i} are
called the faces of the analytic polyhedron. The intersec-
tion of any k different faces (2",;;;;k";;;'n) is said to be an
edge of the analytic polyhedron. If m;;:'n and all faces
have dimension 2n - I, while no edge has dimension
exceeding 2n - k, the analytic polyhedron is a Weil
domain. The set of n-dimensional edges
n ... n
0';, "';, =0';, 0';, forms the skeleton of the
analytic polyhedron. The concept of an analytic
polyhedron is important in problems of integral
representations of analytic functions of several vari-
ables.
References
[I] SHABAT, B.V.: Introduction to complex analysis, 2, Moscow,
1976 (in Russian).
E.D. Solomentsev
Editorial comments. The analytic polyhedron II defined
above is sometimes said to be an analytic polyhedron of
order m (cf. [A1]).
References
[A 1] HORMANDER, L.: An introduction to complex analysis in
several variables, North-Holland, 1973.
AMS 1980 Subject Classification: 32A07
ANALYTIC REPRESENTATION, holomorphic
representation - A representation of a complex Lie
group G in a topological vector space E in which all
matrix elements (cp(g)~, 1/), ~EE, 1/EE', where E' is the
dual topological vector space, are holomorphic on G. A
representation <p is called an anti-analytic representation
if its matrix elements become holomorphic after com-
plex conjugation. An analytic (anti-analytic) representa-
tion of a connected Lie group is uniquely determined
by a corresponding Lie algebra representation of this
group (cf. Representation of a Lie algebra). If G is a
semi-simple complex Lie group, then all its topologi-
cally irreducible analytic (anti-analytic) representations
are finite-dimensional.
References
[I] NAiMARK, M.A.: Theory of group representations, Springer,
1982 (translated from the Russian).
[2] ZHELOBENKO, D.P.: Compact Lie groups and their representa-
tion, Amer. Math. Soc., 1973 (translated from the Russian).
D.P. Zhelobenko
AMS 1980 Subject Classification: 22E10
ANALYTIC RING - A ring of germs of analytic func-
tions at a point in an analytic space. The following is a
190
more accurate definition. Let k be a field with a non-
trivial norm (cf. Nonn on a field) (which is usually
assumed to be complete), and let k { { X I, . . . ,Xn }} be
the k-algebra of power series in XI, ... ,Xn with coef-
ficients in k which converge on some polycylinder with
centre (0, ... ,0), each series converging on its own
polycylinder. A quotient ring of the ring
k {{X I, . . . ,Xn }} is called an analytic ring over k, or
an analytic k-algebra; usually, k is the field of real
numbers R or the field of complex numbers C. Any
analytic ring is a local, Noetherian, Hensel ring; its
field of residues is isomorphic to k. An analytic ring
k {{ X I, . . . ,Xn }} is a regular (and a factorial) ring,
and its completion in the topology defined by the maxi-
mal ideal (X \, ... ,Xn ) coincides with the ring of for-
mal power series k[[X\, ... ,Xn ]]. The normalization
lemma is true: An integral analytic ring is a finite
extension of an analytic ring k { { X \, ... ,Xn }}. Alge-
bras that are finite over k { { X \, ... ,Xn }} are generally
called quasi-analytic k-algebras. If k is a perfect field,
an analytic ring is an excellent ring.
References
[I] DIEUDONNE,1. and GROTHENDIECK, A.: 'Criteres differentiels
de regularite pour les localises des algebres analytiques', J. of
Algebra 5 (1967), 305-324.
[2] MALGRANGE, B.: Ideals of differentiable functions, Tata Inst.
Fundam. Res., 1966.
[3] ABHYANKAR, S.S.: Lacal analytic geometry, Academic Press,
1964.
v.I. Danilov
Editorial comments.
References
[A 1] GRAUERT, H. and REMMERT, R.: Coherent analytic sheaves,
Springer, 1984 (translated from the German).
AMS 1980 Subject Classification: 32CXX, 32810
ANALYTIC SET - A subset of a complete separable
metric space that is a continuous image of the space of
irrational numbers. The concept of an analytic set was
introduced by N.N. Luzin [1]. His classical definition
has been generalized to general metric and topological
spaces.
1) An analytic set in an arbitrary topological space X
is a subset of that space that is the image of a closed
subset of the space of irrational numbers under an
upper semi-continuous multi-valued mapping with com-
pact images of points and a closed graph [2]. If X is a
Hausdorff space, the last-mentioned condition is
automatically satisfied. If X is metrizable, this defini-
tion is equivalent to the classical one.
2) In a complete separable metric space the classical
analytic sets are identical with sf-sets (cf. .sd-set).
This fact forms the base of another definition of ana-
lytic sets (in their capacity as .sd-sets) in general metric
and topological spaces [3], [4], [5]. In the class of
completely-regular spaces, analytic sets in the sense of

l) are absolute analytic sets in the sense of 2). In the
class of non-separable metrizable spaces definition 2) is
employed, since definition l) yields separable analytic
sets.
3) An analytic set in a Hausdorff space [6], [7], is a
continuous image of a subset of a compact space of
type F08.
4) An analytic set is a continuous image of a set
belonging to the family K08, where K is the family of
all closed compact subsets of some topological space
[8]. Definitions l), 3) and 4) equivalent in the class of
Hausdorff spaces.
5) For a generalization in another direction see [4];
k-analytic sets are obtained from closed sets of a topo-
logical space by generalized d -operations (cf. d -
operation; the Baire space of countable weight is
replaced by a Baire space of weight k) and are a gen-
eralization of analytic sets in the sense of 2).
References
[I] LUZIN, N.N.: 'Sur la classification de M. Baire', C.R. Acad.
Sci. Paris Ser. I Math. 164 (1917),91-94.
[2] FROLiK, Z.: 'Respectability of the graphs of composites',
Mathematika 16 (1969), 153-157.
[3] SIERPINSKI, W.: General topology, Univ. Toronto Press, 1956
(translated from the Polish).
[4] STONE, A.H.: 'Non-separable Borel sets II', Gen. Topol. and
Appl. 2 (1972),249-270.
[5] KURATOWSKI, K. and MOSTOWSKI, A.: Set theory, North-
Holland, 1968.
[6] SCHNEIDER, V.E.: 'Descriptive theory of sets in topological
spaces', Uchen. Zap. Moskov. Gos. Univ. 135 (1948), 37-85.
[7] CROQUET, G.: 'Theory of capacities', Ann. Inst. Fourier (Greno-
ble) 5 (1953-1954), 131-295.
[8] SION, M.: 'On analytic sets in topological spaces', Trans. A mer.
Math. Soc. 96 (1960), 341-354.
[9] JAYNE, J.E.: 'Structure of analytic Hausdorff spaces',
Mathematika 23 (1976), 208-211.
[10] ROGERS, c.A., JAYNE, J.E. and DELLACHERIE, c., ET AL. (EDS.):
Analytic sets, Acad. Press, 1980. k such analytic sets are called the irreducible components
A.G. El' in
6) An analytic set in the theory of analytic functions
is a set that can locally be defined as the set of com-
mon zeros of a finite number of holomorphic functions.
If S is an analytic set in an open subset U of the com-
plex n-dimensional space C', this means that for each
point a E U there exist a neighbourhood V C U and a
finite tuple of functions fl' ... ,fr, holomorphic in V,
such that S n V={ZEV: fl(z)= ... =fr(z)=O}. If
the functions J; can be selected (in some neighbour-
hood V) so that the rank of the Jacobi matrix
I 3J; / 3zj I at the point a is r, a is called a regular point
of the analytic set S; the number n - r is called the
(complex) dimension of S at a and is denoted by
<lima s. The set S of all regular points of an analytic
set S is an open everywhere-dense subset of S (in the
induced topology on S as a subset of U). Its comple-
ment S \ S - the set of singular points of S - is an
analytic set in U that is nowhere dense in S.
ANAL YTIC SET
By defmition,
dim" S = linplimz S, a E S;
zeS.
z~a
the dimension of the analytic set S is the number
dim S = sup dim" S.
ae~
An analytic set S is called purely k-dimensional if
<lima S =k for all a ES. For each O~k~dimS, the set
Sk = {a ES: ditna S =k} is a purely k-dimensional ana-
lytic set in U\ U j>kSj' Thus, any analytic set in U
can be represented as a finite union of pure analytic
sets, S = U Sk' At the singular points
<lima(S \ S)<dima s, so that the dimension of the
analytic set of singular points of a purely k-dimensional
analytic set in U is smaller than k. The connected
components of S are complex manifolds. Since this is
also true for the analytic set S \ S, one obtains the
decomposition:
S = S U(S\S)* U ...
of the analytic set into complex manifolds. The decom-
position
S = s; U(S\S;);-l U ...
is more convenient (the dimensions of the summands
strictly diminish, d=dimS); it is called the stratification
of S; the connected components of the k-th summand
of this sum are called the k-dimensional strata of the
analytic set S.
An analytic set S is called reducible (in U) if it is the
union of two analytic sets in U other than itself; other-
wise it is called irreducible (in U). All irreducible ana-
lytic sets in U are connected and pure. An analytic set
S in U is irreducible if and only if the set S of its reg-
ular points is connected. The closure of each connected
component of S is an irreducible analytic set in U;
of S. All analytic sets in U are locally finite unions of
their irreducible components. If two analytic sets have
no common irreducible components, the dimension of
their intersection is strictly smaller than the dimension
of each set. If the intersection of two irreducible sets in
U contains a set that is open in each one, these analytic
sets are identical (the identity theorem).
An analytic set S in U is called irreducible at a point
a ES if there exists a fundamental system of neighbour-
hoods Vi of the point a in U such that all analytic sets
S n Vi in Vi are irreducible; a is then called an irredu-
cibility point of the analytic set S. In a neighbourhood
of each irreducibility point the analytic set has the
structure of an analytic covering, i.e. for each such point
a ES, dima S =k, there exist a connected neighbour-
hood V C U, a linear mapping A: C' ~Ck and an ana-
lytic set a CA(V) such that the restriction of A to S Vn
is a proper mapping into A(V), while the restriction of
A to (S n V)\A-I(a) is a finite-to-one locally biholo-
191
ANALYTIC SET
morphic covering over A(V) \ (1. For irreducible one-
dimensional analytic sets there exists thus (after a suit-
able linear change in coordinates) a local parametric
representation of the form
ZI = ~m, Z2 = /z<n, ... , Zn = f,,<n,
where tEe, I t I <r, m is a positive integer and the
functions fi are holomorphic in the disc I t I <r. Thus,
in a neighbourhood of each irreducibility point, a one-
dimensional analytic set is a topological manifold. For
an analytic set of higher dimension this is usually not
true.
The union of a finite number and the intersection of
any family of analytic sets in V are analytic sets in U.
Any analytic set in V is closed in V. Any compact ana-
lytic set in V cC' consists of a finite number of points.
If V is connected and the analytic set S =1= V, then
V \ S is open, everywhere-dense in V and is also con-
nected. The set of all isolated points of an analytic set
in V has no limit points in V. Moreover, all analytic
sets are locally connected. A connected analytic set is
pathwise connected.
Any analytic set in V of dimension d has locally in V
finite 2d-dimensional Hausdorff measure mes2d. If
diII1a S =k, there exist positive constants c and C
(which depend on a and S) such that
cr 2k ,.;;; meSzk(S n {I Z -a I <r}) ,.;;; Cr 2k,
for all sufficiently small r<O.
The family of analytic sets is invariant under biholo-
morphic mappings (cf. Biholomorphic mapping). More-
over, if S is an analytic set in V and if f: V-,?V I is a
proper holomorphic mapping, then f (S) is an analytic
set in VI'
The definition of analytic sets on complex manifolds
is similar to the definition for C', and all the properties
listed above are preserved except for one: in the general
case there exist compact non-discrete analytic sets. In
special manifolds analytic sets can have certain addi-
tional properties. For example, in the complex n-
dimensional projective space all analytic sets are alge-
braic, i.e. coincide with the set of common zeros of
some finite tuple of homogeneous polynomials (Chow's
theorem).
Real-analytic sets in open subsets of Rn are defined
in the same manner, except that real-analytic functions
must be taken instead of holomorphic functions. Any
real-analytic set is the intersection of some analytic set
(in some open subset of en) with the real subspace
Rn cC'.
References
[1] GUNNING, R.C. and ROSSI, H.: Analyticfunctions of several
complex variables, Prentice-Han, 1965.
[2] HERVE, M.: Several complex variables: local theory, Oxford
Univ. Press, 1963.
E.M. Chirka
192
AMS 1980 Subject Classification: 54H05, 32810,
32C25
ANALYTIC SHEAF - A sheaf F on an analytic space
X such that for any point x EX the set Fx is a module
over the ring (1J x of germs of holomorphic functions at
the point x, and such that the mapping if, a)-'?fa,
defined on the set of pairs (j, a) where fE (1J x' aEFx ,
is a continuous mapping of (1J XF into F for x EX.
M.I. Vnttsekhovskit
Editorial comments.
References
[A1] GRAUERT, H. and REMMERT, R.: Theory of Stein spaces,
Springer, 1979 (translated from the German).
[A2] GRAUERT, H. and REMMERT, R.: Coherent analytic sheaves,
Springer, 1984 (translated from the German)
AMS 1980 Subject Classification: 32L 10
ANALYTIC SPACE - A generalization of the concept
of an analytic manifold. A local model (and, at the
same time, the most important example) of an analytic
space over a complete non-discretely normed field k is
an analytic set X in a domain V of the n-dimensional
space kn over k, defined by equations
fl = ... = /p =0, where fi are analytic functions in V,
which is provided with the sheaf (1J obtained by res-
tricting the sheaf (!) u / I on X; here (!) u is the sheaf of
germs of analytic functions in V, while I is the subsheaf
of ideals generated by fl' ... ,/p. An analytic space
over k is a ringed space that is locally isomorphic to a
ringed space (X, (1J) of the above type. If k is the field
of real numbers R, one speaks of real-analytic spaces; if
k is the field of complex numbers C, of complex-
analytiC spaces or simply of complex spaces; if k is a
field of p-adic numbers Qp' of p-adic analytic spaces.
An analytic (holomorphic) mapping of one analytic
space (X, (!) x) into another (Y, @ y) is a morphism
(X, @ X )-'?( Y, @ y) in the sense of the theory of ringed
spaces, i.e. a pair (</>0, <PI), where </>0: X -'? Y is a continu-
ous mapping, while <PI: <Po I @ y-'? @ x is a sheaf
homomorphism. A point x of an analytic space (X, @)
is called simple (or regular, or non-singular) if x has a
neighbourhood over which (X, @) is isomorphic to a
space of the type (V, @ u), where V is a domain in kn.
Otherwise x is known as a Singular point. A space is
called smooth if all of its points are simple. A smooth
analytic space is identical with an analytic manifold.
The dimension dimx X of an analytic space X at a
point x EX is defined as the dimension of the
corresponding analytic set in a local model (cf. Analytic
set). The global dimension is defined by the formula:
dim X = sUP dim, X.
XEX
Let mx be the maximal ideal of the local nng @~

(x EX). The vector space TxCX)=(mx / m;)* over k is
called the tangent space to X at the point x, while
r;(X)=m x / m; is the cotangent space. The number
emdimx X = dim TAX)
is called the tangent dimension or the embedding dimen-
sion at the point x (the last name is connected with the
fact that em dim X is the smallest number n such that
(X, rrJ) is, in a neighbourhood of x, isomorphic to a
local model in kn). One has dimxX~emdimxX, and
the two are equal if and only if x is a simple point. One
also defines the dimension
em dim X = sUP em dimx X mines an analytic space over k. This correspondence
XEX
Each analytic mapping of analytic spaces
cJ>=(<Po, cJ>\): (X, rrJ x)~(Y, rrJ y) defines a linear map-
pmg dcJ>x: TxCX)~Tcp,,(x)(Y), which IS called its
differential at the point x EX.
An analytic space (X, rrJ) is said to be reduced if its
local model in a neighbourhood of an arbitrary point
has the following property: I consists of all germs of
holomorphic functions that vanish on Xc U. In the
case of an algebraically closed field k, this statement is
equivalent to saying that the fibres rrJ x (x EX) of the
sheaf rrJ contain no nilpotent elements. All smooth
spaces are reduced. If (X, rrJ) is reduced, it can be said
that rrJ consists of the germs of certain continuous
functions on X. The sections of the sheaf rrJ on the
reduced space (X, rrJ) are identical with analytic func-
tions on X, i.e. with analytic mappings X~k (cf. Ana-
lytic mapping). For any analytic space (X, rrJ) there
exists a natural sheaf epimorphism red: rrJ ~ rrJ \ (where
(X, rrJ \) is a reduced analytic space), which is called the
reduction. If jEr(X, rrJ) is a section of the sheaf rrJ,
one can speak of the value of j at a point x EX (which
coincides with the value of the analytic function red j at
x). For this reason the algebra rex,
rrJ), even in the
non-reduced case, is often referred to as the algebra of
analytic (holomorphic) functions on (X, rrJ). Sheaves of
rrJ -modules on an analytic space (X, rrJ) are also called
orzall'tic shoaves.
If (X, rrJ) is an analytic space, then each open U ex
defines an open subspace (U, rrJ Iu). On the other
hand, one can introduce the concept of an analytic sub-
space of (X, rrJ), which is necessarily closed. A set
Y C X is called analytic if it is defined by a finite
number of analytic equations in a neighbourhood of
each point x E X. The sheaf of ideals lye rrJ consisting
of the germs of all analytic functions that vanish on Y,
is connected wi~h such a set. Conversely, each analytic
sheaf of ideals of finite type Ie rrJ defines an analytic
set Y C X. If rrJ y = rrJ / I I y, one obtains an analytic
space (Y, rrJ y), which is called an analytiC subspace of
(X, rrJ); there exists a natural morphism
ANALYTIC SPACE
(1, cJ>\): (Y, rrJ y )~(X, rrJ). An example of an analytic
subspace of (X, rrJ) is the reduction of this space.
The concept of an analytic space originated as a gen-
eralization of the concept of an analytic manifold. Such
a generalization had been suggested mainly by alge-
braic geometry, in which spaces with singular points
had been under study for a long time. The effect of the
ideas of algebraic geometry was immediately reflected
in the ultimate formulation of the concept of an ana-
lytic space (for reduced complex spaces see [9]; for the
general case, see [6]). In particular, any scheme of finite
type over a complete normed field k naturally deter-
between schemes and analytic spaces over k for
reduced complex spaces was studied in [9], in which the
theory of analytic spaces was named 'analytic
geometry'. Subsequently, the two geometries developed
in parallel, and the exchange of ideas between the two
made a substantial contribution to the results achieved
in both these fields.
In the theory of functions of several complex vari-
ables, spaces with singular points appeared, in the first
place, as Riemannian domains (cf. Riemannian
domain), which are analogues of Riemann surfaces for
functions of one variable. Using these as local models,
H. Behnke and K. Stein in 1951 defined a class of
ringed spaces which, as was shown in [5], coincides
with the class of reduced normal analytic spaces (cf.
Normal analytic space). The local geometry of analytic
sets inen had been studied by W. RUckert as early as
1932. Finally, non-smooth analytic spaces are a natural
product of the theory of automorphic functions as quo-
tient spaces of analytic manifolds by properly discrete
groups of automorphisms (cf. Discrete group of
transformations). p-adic analytic sets were first intro-
duced in 1935 by I. Skolem in connection with certain
problems in number theory.
The theory of analytic spaces has two aspects: the
local and the global aspect. Local analytic geometry is
concerned with germs of analytic sets in the space P
provided with sheaves of the above type. Principal
stress is laid on the study of the properties of the alge-
bra of convergent power series in n variables over k
and its quotient algebras - the so-called analytic alge-
bras, the foundations of which were laid by K. Weier-
strass. The local theory comprises the theory of normal-
ization, the study of singular points, local properties of
analytic functions and mappings, etc. The most impor-
tant results obtained in this field refer to the case of
algebraically closed fields k [1], [4], [7]. There appears
the important concept of a coherent analytic sheaf,
which continues to play a leading part in the global
theory. In particular, the structure sheaf rrJ of the ana-
lytic space (X, rrJ) and the sheaf of ideals ly of any
193
ANALYTIC SPACE
analytic set Y e X are coherent for any algebraically
closed k. The case k = R has also been thoroughly stu-
died.
Global analytic geometry studies the properties of
analytic functions, mappings and other analytic objects,
defined 'globally' on the entire analytic space, as well
as the geometrical properties of these spaces. In the
process of studying complex-analytic spaces natural
classes of them were isolated. These include, first, the
class of Stein spaces (cf. Stein space), which can be
roughly described as the class of spaces with a suffi-
ciently large amount of global holomorphic functions.
Stein spaces are the most natural multi-dimensional
generalizations of the domains of the complex plane
considered in the classical theory of functions of one
complex variable. This class of spaces in fact coincides
with the class of analytic subspaces of the spaces en.
Its algebraic analogue is the class of affine algebraic
varieties (cf. Affine variety).
For a domain D ecn holomorphic completeness is
equivalent with the fact that D is a domain of holomor-
phy, i.e. that there exists a holomorphic function in D
that does not extend into a larger domain. The boun-
dary of a domain of holomorphy is pseudo-convex, i.e.
it behaves with respect to local analytic sub manifolds
as would a convex surface with respect to real linear
submanifolds. The problem of the validity of the con-
verse theorem (cf. Levi problem) gave rise to a number
of investigations and yielded a new characterization of
Stein spaces.
The class of compact spaces is, in a certain sense, the
opposite case. The following generalization of the clas-
sical theorem of Liouville is valid: Functions which are
holomorphic on a reduced compact space are constant
on each connected component of this space and
therefore form a finite-dimensional vector space. A gen-
eralization of this theorem are the finiteness theorems,
which confirm the finite dimensionality of the homol-
ogy groups with values in a coherent analytic sheaf.
Holomorphically-convex complex spaces, q-complete,
q-pseudo-convex, q-pseudo-concave spaces, which are
generalizations of Stein spaces, and compact spaces are
also considered (cf. Holomorphically-convex complex
space).
These classes of complex spaces have their analogues
in the theory of holomorphic mappings. Thus, to com-
pact spaces correspond proper holomorphic mappings;
to holomorphically complete spaces correspond Stein
mappings, etc. 'Relative' analogues were found for
many theorems, and the 'absolute' variant of a theorem
is obtained from its relative variant if the entire space
is mapped into a point. The corresponding generaliza-
tion of finiteness theorems are theorems of coherence of
direct images of coherent analytic sheaves under holo-
194
morphic mappings, the first and most important one of
which (for proper mappings) was demonstrated by H.
Grauert [6A).
An important role in the theory of complex spaces is
played by holomorphic mappings of a special kind -
the so-called modifications (cf. Modification), i.e.
mappings f: X ~ Y inducing an isomorphism of open
subspaces X\XI~Y\ YJ, where XI ex, YI eY are
certain analytic sets. One says that Y is obtained from
X by 'contracting' the subset XI on YJ, while X is
obtained from Y by 'blowing up' the subset YI into
X I. Of special interest are analytic subsets that can be
contracted into a point (exceptional analytiC sets); these
were characterized by H. Grauert [6B). A natural prob-
lem in analytic geometry is the following problem of
resolution of singularities: Is it possible to 'blow up' an
analytic space so that the entire space becomes smooth?
It should be noted that modifications in algebraic
geometry were studied as early as the 19-th century,
while modifications in analytic geometry were intro-
duced by Behnke and Stein in 1951 in the context of
the concept of a Riemannian domain.
Another natural object of study, which is also closely
connected with the ideas of algebraic geometry, are
meromorphic functions on complex spaces and their
generalizations - meromorphic mappings (a mapping
which yields an operation inverse to a modification
may serve as an example; cf. Meromorphic function;
Meromorphic mapping). Meromorphic functions on a
reduced compact complex space X form a field of tran-
scendence degree t(X),;;;;dimX (this was first demon-
strated by c.L. Siegel in 1955 for the smooth case).
Spaces X for which t(X)=dimX (Moishezon spaces)
form a class which is very close to the class of projec-
tive algebraic varieties; they are characterized by the
fact that they are modifications of smooth projective
algebraic varieties. Another class of analytic spaces
which is very close to algebraic varieties, are Kahler
manifolds (cf. Kiihler manifold). A number of criteria
for the projectivity of a compact complex space are
known [3], [6B], [13]. Studies of automorphic functions
in several complex variables have made a major contri-
bution to the development of this subject.
The theory of deformations of analytic structures (cf.
Deformation) is concerned with the problem of classifi-
cation of analytic objects of a given type (e.g. all com-
plex structures on a given real-analytic variety, all ana-
lytic subspaces in a given complex space, etc.), with the
purpose of introducing the 'natural' structure of a com-
plex space on the set of these objects, and in order to
describe all analytic objects 'sufficiently near' to the
given object. In the former case one speaks of the glo-
bal moduli problem, while in the latter one speaks of
the local moduli problem. An example of the global

moduli problem is the problem of the classification of
all complex structures on a compact Riemann surface
(cf. Moduli of a Riemann surface).
The principal apparatus of global analytic geometry
is formed by coherent analytic sheaves and their coho-
mology spaces. The first successful result of the coho-
mological method was Cartan's solution of the additive
Cousin problem and of problems of prolongation of a
holomorphic function from a closed Y C X for Stein
manifolds (X, l!) (cf. Cousin problems; [8]); it was
found that the solution of these problems is obstructed
by the cohomology groups HJ(X, l!) and HJ(X, ly),
respectively.
Most of the results of the global theory were initially
demonstrated for complex manifolds, after which they
were generalized to analytic spaces. The difficulties
involved in this procedure often necessitated the
development of completely new methods. Cohomology
spaces of a locally free analytic sheaf on a complex
manifold may be expressed in terms of differential
forms (the Dolbeault - Serre theorem, cf. also Differen-
tial fonn), which makes it possible to study them by
methods of the theory of elliptic differential equations
and by other analytic methods. In the non-smooth case
this approach involves major difficulties, and it is often
necessary to define cohomology classes in other ways,
e.g. using Cech cochains in a suitable covering. The
technique of Banach analytic spaces, applied to moduli
problems, proved useful in this context (d. Banach ana-
lytic space).
See also Real-analytic space; Rigid analytic space.
References
[1] ABHYANKAR, S.S.: Local analytic geometry, Academic Press,
1964.
[2] BANICA, C. and STANASILA, 0.: Algebraic methods in the global
theory of complex spaces, Wiley, 1976 (translated from the
Rumanian).
[3] GUNNING, R.C. and ROSSI, H.: Analytic functions of several
complex variables, Prentice-Hall, 1965.
[4] GRAUERT, H. and REMMERT, R.: Analytische Stellenalgebren,
Springer, 1971.
[5] GRAUERT, H. and REMMERT, R.: 'Komplexe Rliume', Math
Ann. 136 (1958),245-318.
[6A] GRAUERT, H.: 'Ein Theorem der analytischen Garbentheorie
und die Modulrliume komplexer Strukturen', Publ. Math
IHES 5 (1960), 233-292.
[6B] GRAUERT, H.: 'Ueber Modifikationen und exzeptionelle
analytische Mengen', Math Ann. 146, no. 4 (1962), 331-368.
[7] NARASIMHAN, R.: Introduction to the theory of analytic spaces,
Springer, 1966.
[8] CARTAN, H.: 'Varietes analytiques complexes et cohomologie',
in Colloque sur les fonctions de plusieurs variables. Bruxelles, 11
au 14 mars 1953, G. Thone and Masson, Liege-Paris, 1953.
[9] SERRE, J.-P.: 'Geometrie algebrique et geometrie analytique',
Ann. Inst. Fourier (Grenoble) 6 (1955-1956), 1-42.
[10] FUKS, B.A.: Special chapters in the theory of analytic functions
of several complex variables, Amer. Math. Soc., 1965 (translated
from the Russian).
[II] HORMANDER, L.: An introduction to complex analysis in several
variables, North-Holland, 1973.
ANALYTIC SURFACE
[12] lIIRzEBRUCH, F.: TopolOgical methods in algebraic geometry,
Springer, 1978 (translated from the German).
[13] CHERN, S.S.: Complex manifolds without potential theory,
Springer, 1979.
A.L. Onishchik
Editorial comments. Instead of [8] one may also consult
[A2].
The original papers in which what are now called
Moishezon spaces were introduced are [A3]. They are a
natural setting to examine the question of projectivity of
compact complex spaces. In particular, part I contains the
result that a Moishezon space is algebraiC projective if and
only if it carries a Kahler metric.
References
[A 1] GRAUERT, H. and REMMERT, R.: Coherent analytic sheaves.
Springer. 1984 (translated from the German).
[A2] CARTAN. H.: Sem. ENS. 1951-1952. Ecole Norm. Super ..
[A3] MOISHEZON, B.: 'On n-dimensional compact complex
varieties with n algebraically independent meromorphic func-
tions'. Amer. Math. Soc. Translations Ser. 263 (1967). 51-
177.
AMS 1980 Subject Classification: 32CXX
ANALYfIC SURFACE in a Euclidean space - An
arbitrary two-dimensional analytic submanifold X in
the space Rn , n > 2. However, the term 'analytic surface
in Rn ' is often employed in a wider sense as a manifold
which is (locally) analytically parametrizable. This
means that the coordinates of the points
X=(Xb'" ,Xn)EX can be represented by analytic
functions Xj =Xj(t) of a real parameter t =(tj, ... ,tk)
which varies in a certain range aCRk , l.;;;;k<n. If the
rank of the Jacobi matrix II ax / at II, which for an
analytic manifold is maximal everywhere in a, is equal
to k, then the dimension of the analytic surface X is k.
In the complex space en the term 'analytic surface' is
also employed to denote a complex-analytic surface X in
en, i.e. a manifold which allows a holomorphic
(complex-analytic) parametrization. This means that
the complex coordinates of points z=(zJ, ... ,Zn)EX
can be expressed by holomorphic functions Zj =Zj(T) of
a parameter T = (TJ, . . . ,Tk) which varies within a cer-
tain range acck , l.;;;;k<n (it is also usually assumed
that rank II az / aT II =k). If a =c" and all the func-
tions Zj(T) are linear, one obtains a complex-analytic
plane (cf. Analytic plane). If k = 1, the term which is
sometimes employed is holomorphic curve (complex-
analytic curve); if all functions Zj(T) are linear, one
speaks of a complex straight line in the parametric
representation:
Zi = aiT+bi; ai,biEC, i=l, ... ,n, (al> ... ,an) =1= O.
References
[I] SHABAT, B.V.: Introduction to complex analySiS, 1-2, Moscow,
1976 (in Russian).
[2] VLADIMIROV, V.S.: Methods of the theory offunctions of several
195
ANALYTIC SURFACE
complex variables, M.I.T, 1966, Chapt. 2 (translated from the
Russian).
E.D. Solomentsev
E.M. Chirka
AMS 1980 Subject Classification: 32C25, 32B15
ANALYTIC SURFACE (IN ALGEBRAIC
GEOMETRY) - A two-dimensional (complex) analytic
manifold, i.e. a smooth four-dimensional manifold with
a complex structure. While the theory of analytic sur-
faces forms part of the general theory of complex mani-
folds, the two-dimensional case is treated separately,
since much more is known about analytic surfaces than
about n-dimensional manifolds if n ;;;;'3. Moreover, cer-
tain facts are specific to the two-dimensional case
alone. These results concern the classification of ana-
lytic surfaces, which is analogous to that of algebraic
surfaces (cf. Algebraic surface) - a fact which largely
reduces the theory of analytic surfaces to that of alge-
braic surfaces. The principal results on the classification
of analytic surfaces were obtained by K. Kodaira [l],
[2], [3], but his work is based on the results of the clas-
sical Italian school of algebraic geometry on the classif-
ication of algebraic surfaces.
All analytic surfaces discussed below are assumed to
be compact and connected.
Examples. I) Algebraic surfaces. Let
j;(xo, ... ,XN), i=l, ... ,m,
be a set of homogeneous polynomials with complex
coefficients. The closed subset of the complex projective
space pN (C) specified by the equations ii(x) = 0 is an
analytic surface if it is non-singular, connected and has
complex dimension two. This is the basic example of an
analytic surface.
2) Complex tori. Let C2 be the two-dimensional vec-
tor space over the field of complex numbers (as a vec-
tor space over the field of real numbers it is isomorphic
to R4) and let r"-'z4 be a lattice in C2. The quotient
space X = C2 / r is an analytic surface. Being a smooth
manifold, X is diffeomorphic to a four-dimensional
torus, but the complex structure on X depends on the
lattice r. Complex tori X = C2 / r play an important
role in analysis, since meromorphic functions on such
tori are meromorphic functions on C2 and are periodic
with period lattice r. Analytic surfaces of the type
C2 / r are not always algebraic. There also exist lattices
r such that there are no meromorphic functions at all
(except for constants) on the corresponding torus
C2 / r. For specific examples of such tori see [5].
3) Hopf surfaces. Let Y=C 2 -{0} and let c be a
positive number. Consider the action of the group Z on
Y given by
196
The group Z acts discretely and without fixed points on
Y, while the quotient space X = Y / Z is diffeomorphic
to SI XS3. The quotient space X has a natural struc-
ture of an analytic surface, and is called a Hopf surface.
Oassification of analytic surfaces. The principal
invariant in the classification of analytic surfaces is the
transcendence degree of the field of meromorphic func-
tions C(X) on the analytic surface X. According to
Siegel'S theorem, for any compact connected manifold
X the field C(X) is finitely generated, and its transcen-
dence degree is not larger than the complex dimension
of X. Thus, for an analytic surface X, the field C(X)
contains two independent meromorphic algebraic func-
tions, or one such function or constants only. These
possibilities lead to the following theorems.
For any analytic surface X to be an algebraic surface
it is necessary and sufficient that there exist two alge-
braically independent meromorphic functions on X.
If an analytic surface X has a field of meromorphic
functions of transcendence degree I, then X is an ellip-
tic surface, i.e. there is a holomorphic mapping onto an
algebraic curve Y, P : X ~ Y, such that
P'C(Y) = C(X)
and all fibres of P, except for a finite number, are ellip-
tic curves (the singular fibres may only have a very spe-
cial form, which has been thoroughly studied).
If no meromorphic functions other than constants
exist on an analytic surface X, and no exceptional
curves (cf. Exceptional subvariety) exist on X, then the
first Betti number b I of X assumes only three values:
4, I or O. If b I = 4, X is a complex torus, and if b I = 0,
X has a trivial canonical fibration. These analytic sur-
faces are called K3-surfaces. They are all mutually
homeomorphic. The case b I = 1 has not been studied in
detail, but certain examples of analytic surfaces with
b I = I are obtained by a generalization of the construc-
tion of Hopf surfaces.
Analytic Kahler surfaces are not always algebraic.
They are, however, algebraic if the square of their first
Chern class is positive. All analytic Kahler surfaces
with b I >0 are deformations of algebraic surfaces.
References
[I] KODAIRA, K.: Matematika 6, no. 6 (1962), 3-17.
[2A] KODAIRA, K.: 'On compact (complex) analytic surfaces, 1',
Ann. of Math. 71, no. 1(1960), 111-152.
[2B] KODAIRA, K.: 'On compact (complex) analytic surfaces. II',
[2C]
Ann. of Math. 77, no. 3 (1963), 563-626.
KODAIRA, K.: 'On compact (complex) analytic surfaces, Ill',
Ann. of Math. 78, no. 1(1963), 1-40.
[3A] KODAIRA, K.: 'On the structure of compact (complex) analytic
surfaces 1', Amer. 1. Math. 86 (1964), 751-798.
[3B] KODAIRA, K.: 'On the structure of compact (complex) analytic
surfaces II', A mer. 1. Math. 88 (1966),682-721.
[3C] KODAIRA, K.: 'On the structure of compact (complex) analytic
surfaces III', A mer. 1. Math. 90 (1968), 55-83.
[3D] KODAIRA, K.: 'On the structure of compact (complex) analytic
 ANALYTIC THEORY OF DIFFERENTIAL EQUATlO'JS

surfaces IV', Amer. J. Math. 90 (1968), 1048-1066.
[4) 'Algebraic surfaces', Trudy Mat. Inst. Steklov. 75 (1965) (in
Russian).
[5) SHAFAREVICH, I.R.: Basic algebraic geometry, Springer, 1977
(translated from the Russian).
B.B. Venkov
Editorial comments. The notion defined above is also
called a complex-analytic surface, since one considers com-
plex structures and the field of complex numbers C. If
instead one considers real structures and the field of real
numbers, one speaks of real-analytic surfaces. However, an
analytic surface is always understood in the sense explained
above.
References
[A1] BARTH, W., PETERS, C. and VEN, A. VAN DER: Compact com-
plex surfaces, Springer, 1984.
AMS 1980 Subject Classification: 14JXX
ANALYTIC THEORY OF DIFFERENTIAL EQUA-
TIONS - The branch of the theory of ordinary differen-
tial equations in which the solutions are studied from
the point of view of the theory of analytic functions. A
typical formulation of a problem in the analytic theory
of differential equations is this: Given a certain class of
differential equations, the solutions of which are all
analytic functions of one variable, find the specific pro-
perties of the analytic functions that are solutions of
this class of equations. In this wide sense, the analytic
where A _j are constant matrices, A -P7'~:0, and the
matrix B(t) is holomorphic at to. A pole 10*,00 of
order p is called a regular Singular point if p = 1 and an
irregular singular poinl if p~2. The case 10 = 00 is
reduced to the case 10 = 0 by the change of variables
I~/-I. In what follows, 10*,00.
2) Let 10 be a pole of A (I). Then there exists a fun-
damental matrix X(/) for the system (2) of the form
X(t) = 1I>(t)(t-tof, (3)
where D is a constant matrix, 4>(/) is holomorphic if
I t - 10 I < p where to is a regular singular point, and
4>(/) is holomorphic if 0< I t - to I <p where to is an
irregular singular point, for some p>O. (Here,
(I - (0)D = exp(ln(t - (0)D), by definition.) For a regular
singular point the matrix D can be expressed in terms
of A (t) in an explicit form [1], [2]; this is not the case
for irregular singular points.
A similar classification of singular points is intro-
duced for differential equations of order n with mero-
morphic coefficients. Differential equations and dif-
ferential systems all singular points of which are regular
are known as Fuchsian differential equations (systems).
The general form of A (t) for such a system is:
k
A(t) = ~(t-tj)-IAj, Aj = const, k<oo.
j=l

theory of differential equations includes the theory of An example of a Fuchsian differential equation is the
algebraic functions, the theory of Abelian integrals, the hypergeometric equation.
theory of special functions, etc. Special functions - 3) Let A (t)=tqB(t) where q~O is an integer and let
Bessel functions, Airy function", Legendre functions, B(/) be holomorphic at 1=00 (00 is an irregular singu-
Laguerre functions, Hermite functions (cf. Hermite lar point if B( 00 )*,0). If S is a sufficiently narrow sec-
function), Chebyshev functions (cf. Chebyshev function), tor of the form I t I >R, a<argt<{3, then there exists
Whittaker functions, Weber functions (cf. Weber func- a fundamental matrix of the form
tion), Mathieu functions, hypergeometric functions (cf. X(t) = P(t)tQexp(R(t, (4)
Hypergeometric function), Sonin functions and many
other functions - are solutions of linear differential where Q is a constant matrix, R (t) is a diagonal matrix
equations with analytic coefficients. whose elements are polynomials in t 1/ p, P ~ 1 is an
integer and
Linear theory. Consider a system of n equations in 00
pet) ~ L,Pjt- Jlp
matrix notation: j=O
x= A (t)x +f(t). as I t I ~OO, tES. The plane C(/) is subdivided into a
(I)
1) Let the matrices A (t),f(t) be holomorphic in a finite number of sectors, and in of them there exists a
region G cC(t), where C(t) is the complex t-plane. Any fundamental matrix of the form (4) ([3], [4]; see also [1],
solution of the system (1) will then be analytic in G [2]).
(but will not, in general, be single-valued if G is not 4) As a result of analytic continuation along a closed
simply-connected). It is assumed that A (t) is mero- path y the fundamental matrix X(/) is multiplied by
morphic in G, and one considers the homogeneous sys- By: X(t)~X(t)By, where By is a constant matrix; one
tem obtains the monodromy group of the differential equa-
x = A (t)x. (2)
tion. LA. Lappo-Danilevskii [5] has studied the prob-
(The matrix A (t) is called holomorphic (meromorphic) in lem of Riemann: Let A (t) be a rational function of t
G if all its elements are holomorphic (meromorphic) in and let the singularities of the fundamental matrix X(t)
G.) A point to EG is called a pole of the matrix A (t) of be known, find A (I).
order p~ 1 if, in a given neighbourhood of this point, 5) Let the function z =CP(/) be a conformal mapping
A(t) = A_.(t-to)-+ +A_1(t-to)-I+B(t), of the upper half-plane Imt>O onto the interior of a

197
ANALYTIC THEORY OF DIFFERENTIAL EQUATIONS
polygon, the boundary of which consists of a finite
number of segments of straight lines and circular arcs.
The function cf>(t) will then satisfy the Schwarz equa-
tion: Z '" 3 ["
z ]2
{z, t} i-"2 i = R(t).
where R (t) is a rational function, and the equation
" 1
w +"2R(t)w =0
is Fuchsian. Any solution of equation (5) may be
represented in the form z = W I / W 2, where W I and W 2
are linearly independent solutions of equation (6). Let
G be an infinite discrete group and let cf>(t) be an auto-
morphic function of G, then cf>(t) can be represented as
</>=WI/W2' where WI, W2 are linearly independent
solutions of equation (6) and R (t) is some algebraic
function.
Non-linear theory. 1) Consider the Cauchy problem:
x = f(t, x), x(to) = xo,
where tEC(t), x=(xJ, ... ,Xn)Ecn(X),
j=(jJ, ... ,J,,).
Cauchy's theorem: Let the function j(t, x) be holo-
morphic in (t, x) in a region GCC(t)xcn(x) and let
the point (to, XO)EG. Then there exists a 8>0 such
that in the domain I t - to I <8 there exists a solution
x(t; to, xo) of the Cauchy problem (7), which is unique
and holomorphic.
An analytic continuation of the solution x(t;to,xo)
will also be a solution of the system (7), but the func-
tion obtained as a result of the continuation may have
singularities and, in the general case, is a many-valued
function of t. The problems which arise are: What
singularities may this function have and how can one
construct the general solution? In the linear case these
questions have been conclusively answered. In the
non-linear case the situation is much more complicated
and has not been fully clarified even when the !jet, x)
are rational functions of t, x.
2) Consider the differential equation:
dx_~
dt - Q(t, x)'
where t E C, X E C, and P and Q are holomorphic func-
tions of (t, x) in a certain region G. A point (to, xo) is
called an (essentially) Singular point of equation (8) if
P(to, xo)=O, Q(to, xo)=O. Below the structure of the
solutions in a neighbourhood of a singular point of the
equation is clarified. Develop P and Q into Taylor
senes:
P(t,x) = all(x-xO)+a12(t-tO)+ ... ,
Q(t,x) = a2l(x-xo)+adt-t o)+ : .. ,
and let AI, A2 be the eigen values of the matrix I au II
The following theorem holds. Let Aj=l=O and let none of
198
the numbers Al / A2, A2 / Al be either a non-negative
integer or a real negative number. Then there exists a
neighbourhood U of the p.,oint (to, xo), a neighbour-
hood V of the point t =0, =0, and functions x
(5) (=(t, x) and x=x(t, x) such that the mapping U~V
defined by these functions is biholomorphic, and the
differential equation (8) in the new variables assumes
the form [6]:
(6) dx AIX
dt A2t
All solutions of equation (8) in the new variables are
written in the form x=C?,,/lo., and x=O. Thus, a
singular point of the equation is a branching point of
infinite order for all solutions of equation (8) (except
for the trivial solutions). The singular points of the
solution which coincide with the singular points of the
equation are called stationary. As distinct from the
linear case, the solution of a non-linear equation may
have singular points not only at the singular points of
(7)
the equation; such singular points of the solution are
called movable. Painleve's theorem is valid: The solu-
tions of the equation
P(t, x, x) = 0,
where P is a polynomial in x and x with holomorphic
coefficients in t, has no movable transcendental singu-
lar points [7].
If, in equation (8), P and Q are polynomials in t, x,
then, in view of Painleve's theorem, all movable singu-
lar points are algebraic. On substituting t = 1 / t',
x =x' / t', equation (8) assumes the form
dx' PI(t', x')
dt' = -Q-'-I-'-(t7"",,-x-:'-,),
where P I and QI are polynomials. Let Xj be the roots
of the equation PI (0, x') =0. The points (0, xj) are
called infinitely-remote Singular points of equation (8);
the structure of the solutions in a neighbourhood of
these points is described by the theorem quoted above
[6].
Let P and Q be polynomials of degree n. Since P, Q
(8) are defined by their coefficients and the pair (AP, AQ)
defines the same equation, one obtains a one-to-one
correspondence between equations (8) and the points of
the complex projective space CpN,
N=(n + l)(n -2)-1. The following theorem is valid: If
some set of measure zero is removed from CpN, the
remaining equations (8) will have the following pro-
perty: All solutions x=x(t) are everywhere dense in
C2 =C(t)XC(x) [8].
3) Consider the autonomous system
x = f(x), (9)
tEC(t),xEcn(X),j=(j], ... ,J,,). A point xU will
then be a singular point oj the system (9) if j(xo) = o.

Poincare's theorem is valid: Let xO be a singular point
of the autonomous system (9). Also let a) the elemen-
tary divisors of the Jacobi matrix / (xo) be prime divi-
sors; and b) the eigen values Aj of this matrix lie on
one side of some straight line in C(A) passing through
the coordinate origin. Then there exists neighbourhoods
U, V of the points x = x ,x
= 0 and a biholomorphic
mapping x =x(x): V ~U such that the system (9)
x
expressed in the variable assumes the form [9]:
iXJ - G with values in V (cf. Representation theory). If T is a
dt = AjXj, 1~j ~n.
If only condition a) is satisfied, it is possible, by using
a transformation x =cp(x), where q,(x) is a formal power
series, to convert system (9) in a neighbourhood of a
singular point into a system which can be integrated in
quadratures [9], [10]. However, the convergence of these
series has been proved on assumptions close to a) and
b). If the function f(x) and the transformation x =q,(x)
are real for real x, x,
a theorem similar to the theorem
of Poincare has been proved [11]. The structure of the
solutions of the autonomous system (9) in general,
where jj(x) are polynomials and n;;;' 3, has not yet
(nineteen-seventies) been studied.
References
[1) CODDINGTON, E.A. and LEVINSON, N.: Theory of ordinary dif-
ferential equations, McGraw-Hill, 1955.
(2) WAZOV, W.: Asymptotic expansions for ordinary differential
equations, Interscience, 1965.
[3) BIRKHOFF, G.D.: 'Singular points of ordinary linear differential
equations', Trans. A mer. Math. Soc. 10 (1909),436-470.
(4) 'fIuITZINSKY, W.J.: 'Analytic theory of linear differential equa-
tions', Acta Math. 62 (1934), 167-226.
(5) LAPpo-DANILEVSKY, J.A.: Memoire sur la tMorie des systemes
des equations differentielles lineaires, Chelsea, reprint, 1953.
(6) BIEBERBACH, L.: Theorie der gewohnlichen Differential-
gleichungen auffunktionentheoretischer Grundlage dargestellt,
Springer, 1965.
[7) GOLUBEV, V.V.: Vorlesungen Uber Differentialgleichungen im
Komplexen, Deutsch. Verlag Wissenschaft., 1958 (translated
from the Russian).
(8) KHuoAi-VERENOV, M.G.: 'On a property of solutions of a dif-
ferential equation', Mat. Sb. 56 (98), no. 3 (1962), 301-308 (in
Russian).
(9) NEMYTSKII, V.V. and STEPANOV, V.V.: Quantitative theory of
differential equations, Princeton Univ. Press, 1960 (translated
from the Russian).
(10) BRYUNO, A.D.: Local methods of non-linear analysis of differen-
tial equations, Moscow, 1979 (in Russian).
[11) SIEGEL, c.L.: 'Ueber die Normalform analytischer Differential-
gleichungen in der Nilhe einer Gleichgewichtslosung', Nachri-
chten Akad Wissenschaft. Gottingen (1952), 2130.
(12) POINCARE, H.: Oeuvres de Henri Poincare, Vol. 3, Gauthier-
Villars, 1916-1965.
(13) FORD, L.R.: Automorphicfunctions, Chelsea, reprint, 1951.
M. V. Fedoryuk
Editorial comments. The Riemann monodromy problem
mentioned above is of great importance in the modern
theory of completely integrable or soliton equations. Cf.
Soliton.
A:--:AL YTIC VECTOR
References
[A 1] INCE, E.L.: Ordinary differential equations, Dover, reprint,
1956.
[A2] HILLE, E.: Ordinary differential equations in the complex
domain, Wiley (I nterscience) , 1976.
AMS 1980 Subject Classification: 34A25
ANALYTIC VECTOR in the space V of a representa-
tion T of a Lie group G - A vector ~ E V for which the
mapping g~T(g)~ is a real-analytic vector function on
weakly-continuous representation of a Lie group G in a
Banach space V, then the set VW of analytic vectors is
dense in V [I], [2], [3]. This theorem has been general-
ized to a wide class of representations in locally convex
spaces (5). It has been proved (6) that a representation
of a connected Lie group G in a Banach space V is
uniquely determined by the corresponding representa-
tion of the Lie algebra of the Lie group G. in the space
VW.
An analytic vector for an unbounded operator A on a
Banach space V, defined on a domain D (A), is defined
as a vector 00
~ E nD(An)
n=1
for which the series
has a positive radius of convergence. This notion, which
was introduced in [2], is a special case of the general
concept of an analytic vector; here the set of points on
the real line with the addition operation plays the role
of the Lie group G. It was found useful in the theory of
operators on Banach spaces and in the theory of elliptic
differential operators.
References
[I) CARTIER, P. and DIXMIER, J.: 'Vecteurs analytiques dans les
representations de groupes de Lie', A mer. J. Math. 80 (1958),
131-145.
(2) NELSON, E.: 'Analytical vectors', Ann. of Math. 70 (1969),572-
615.
[3) GARoING, L.: 'Vecteurs analytiques dans les representations
des groups', Bull. Soc. Math. France 88 (1960),73-93.
[4] CARTIER, P.: 'Vecteurs analytiques', in Sem. Bourbaki
1958/1959, Vol. 181, 1959, pp. 12-27.
(5) MOORE, R.T.: 'Measurable, continuous and smooth vectors for
semigroup and group representations', Mem. A mer. Math. Soc.
78 (1968).
[6) HARISH-CHANDRA: 'Representations of a semisimple Lie group
on a Banach space 1', Trans. A mer. Math. Soc. 75 (1953), 185-
243.
A.A. Kirillov
Editorial comments.
References
[A1] WARNER, G.: HarmOniC analysis on semi-simple Lie groups,
1, Springer, 1972.
AMS 1980 Subject Classification: 22E45
199
ANDRONOV - WITT THEOREM

ANDRONOV - Wrrr THEOREM - A modification

of Lyapunov's theorem (on the stability of a periodic
J

(z) = sin /!7T
V7T
[1- _Z_2_
22 -,?-
+ --::---:=-z_4-:---::-_
(22 - '?-)(42 -'?-)
solution of a non-autonomous system of differential
equations) for the autonomous system - (22_'?-)(42~'?-)(62_,?-) + ... ] +
dx,
dt = X;(xJ, ... ,xn), i=I, ... ,n. (1) + sin /!7T Z_ _
[ __ Z3 +
7T 12-,?- (12-'?-)(3 2 -,?-)
Let
= eMt) (2)
(l2-'?-)(32~'?-)(52-,?-)
X,

be a periodic solution of the system (1), and let + ... }

. ~aX'(</>I""'</>n) . The asymptotic expansion

~, = ..::.. a. ~)' 1= 1, ... , n, (3)
j= I x)
Jp(z) ~ J.(z)+
be the corresponding system of variational equations
+ sinll7T [1- 12-,?- + (12-'?-)(3 2 -,?-) _ ... ]_
which always has, in the case here considered, one zero 7TZ Z2 z4
characteristic exponent. The Andronov- Witt theorem
is then valid: If n - 1 characteristic exponents of the _ sinV7T [.!:..._ v(2 2 -,?-) + v(22 -,?-)W-'?-) _ ... ]
7TZ Z Z3 Z5
system (3) have negative real parts, a periodic solution
(2) of the system (1) is stable according to Lyapunov is valid for I z I-HXJ and I argz I <7T.
(cf. Lyapunov characteristic exponent; Lyapunov stabil- The functions have been named after CT. Anger [1],
ity). who studied functions of the type (*), but with 27T as
The Andronov- Witt theorem was first formulated the upper limit of the integraL
by A.A. Andronov and A.A. Witt in 1930 and was References
proved by them in 1933 [1]. [1] ANGER, C.T.: Neueste Schr. d. Naturf d. Ges. i. Danzig 5
(1855), 1-29.
References [2] WATSON, G.N.: A treatise on the theory oj Bessel junctions,
[1] ANDRONOV. A.A.: Collected works, Moscow. 1976 (in Russian).
Cambridge Univ. Press, 1952.
[2] PONTRYAGIN, L.S.: Ordinary differential equations, Addison- A.P. Prudnikov
Wesley, 1962, 264 (translated from the Russian).
AMS 1980 Subject Classification: 33A40
E.A. Leontovich-Andronova
Ed ito rial comments. The Andronov - Witt theorem is ANGLE - A geometrical figure consisting of two dis-
usually found in the Western literature under some heading tinct rays issuing from the same point. The rays are
like' hyperbolic periodic attractor'. called the sides of the angle and their common original
Good additional general references are [A 1], [A2], [A3]. is called the vertex of the angle. Let [BA), [BC) be the
In [A2] the theorem under discussion occurs as a statement
sides of an angle, B its vertex and a the plane defined
about periodic attractors, cf. pp. 277-278. The original
by its sides. The figure r=[BA) U [BC) divides the
Andronov-Witt paper is [A4].
plane a into two adjoint regions a;, i = 1, 2:
References
a'l U a~ =a \ r. The region ai =a; U r, i = 1,2, is also
[A1] HAHN, W.: Stability of motion, Springer, 1967, 422.
[A2] HIRSCH, M. and SMALE, S.: Differential equations, dynamiC called an angle or a plane angle, a; is called the interior
systems and linear algebra, Acad. Press, 1974. domain of the plane angle ai' The regions (angles)
[A3] CODDINGTON. E.A. and LEVINSON, M.: Theory of ordinary
aI, a2 are called supplementary.
differential equations, McGraw-Hili, 1955, 323.
[A4] ANDRONOV. A.A. and WITT, A. 'Zur Stabilitat nach Two angles are called equal (or congruent) if they can
Liapounoy', Physikal. Z. Sowjetunion 4 (1933),606-608. be superposed so that corresponding sides and vertices
AMS 1980 Subject Classification: 34020, 34C25 coincides. For every ray in the plane and a given side
of it one can construct a unique angle equal to a given
ANGER FUNCfION - The function one. The comparison of two angles can be realized in

f
two ways. If the angle is considered as a pair of rays
Jp(z) = 1- cos(v</> - z sin </ d</>, (*) with a common origin then to answer the question
7T 0
which of the two angles is the greater it is necessary to
which satisfies the inhomogeneous Bessel equation: superpose in the given plane the vertices and one of
2 I . each pair of sides (see Fig. I). If the second side of an
rv- "+ZV
7" ,
+(z -,r-)"
7
= -(z-v)smV7T.
. ~ 11
angle comes to lie inside the second angle, one says
For integers v=n, Jv(z)=Jn(z) is the Bessel function of that the first angle is smaller than the second.
order n (cf. Bessel functions). For non-integer v the fol- A second method of comparing angles consists in
lowing expansion is valid: assigning to each angle a certain number. Equal angles

200

Fig. I.
will have the same number of degrees or radians (see
below), greater angles a greater and smaller angles a
smaller number.
Two angles are called supplementary if they have
their vertex and one side is common, while the other
two sides form a straight line (see Fig. 2).
L
---
[ 0 8 tain angle cp - the number given by the vector formula
Fig. 2.
Generally, angles having a common vertex and one
common side are called a4joining. Two angles are called
vertical if the sides of one are the prolongations through
the vertex of the sides of the other. Vertical angles are
equal to each other. Angles the sides of which form a
straight line are called straight. Half a straight angle is
called a right angle. A right angle may equivalently be
defined as follows: An angle equal to its supplement is
called right. The interior domain of a plane angle that
is not greater then a straight angle is called a convex
domain in the plane.
As the unit of measurement of an angle one uses the
90-th part of a right angle, called a degree. The cyclic
or radian measure of an angle is also used. The numeri-
cal value of the radian measure of an angle is the
length of the arc cut out by the sides of the angle on
the unit circle. One radian is assigned to the angle the
arc of which has length equal to the radius. A straight
angle is equal to 1T radians.
By intersecting two straight lines lying in a plane
with a third there are formed angles (see Fig. 3): 1 and
5, 2 and 6, 4 and 8, 3 and 7; these are called
corresponding angles; 2 and 5, 3 and 8 are called interior
angles on the same side; 1 and 6, 4 and 7 are called
exterior angles on the same side; 3 and 5, 2 and 8 are
called interior opposite angles; and I and 7, 4 and 6 are
called exterior opposite angles.
Fig. 3.
ANGULAR BOUNDARY VALUE
In practical problems it is convenient to consider
angles as the measure of rotation of a fixed ray round
its origin to a given position. Depending on the direc-
tion of the rotation, the angle may in this case be con-
sidered as positive or negative. Thus, an angle in this
sense can have as value any real number. The angle as
a measure of the rotation of a ray is considered in the
theory of trigonometric functions: For any value of the
argument (an angle) one can define a value for the tri-
gonometric functions. The concept of an angle in
geometrical systems founded on the axiomatics of posi-
tions of point vectors is radically different from the
definition of angles as figures - in such axiomatics by
an angle one understands a well-defined metric quan-
tity connected with two vectors by means of their scalar
product. Thus, each pair of vectors a, b defines a cer-
_ (a, b)
cos</> - I a 1'1 b I '
where (a, b) is the scalar product of the vectors.
The concept of an angle as a plane figure and as a
certain numerical quantity is applied in different
geometrical problems in which the angle plays a special
role. Thus, by the angle between two intersecting curves
having a well-defined tangent in their point of intersec-
tion one understands the angle formed by these
tangents.
By the angle between a straight line and a plane one
understands the angle formed by the straight line and
its orthogonal projection on the plane; its value lies
between 0 and 90. By the angle between two crossing
lines one understands the angle between the directions
of these lines, i.e. between two lines parallel to the
given crossing lines and passing through a given point.
A solid angle is a part of space bounding a certain
conical surface; a particular case of a solid angle is a
polyhedral angle.
In higher-dimensional geometry one defines the angle
between multi-dimensional planes, between straight
lines and planes, etc. Angles between straight lines,
between planes and straight lines, and between higher-
dimensional planes are also defined in non-Euclidean
spaces.
LA. Sidorov
Editorial comments.
References
[A1] GREENBERG, M.J.: Euclidean and non-Euclidean geometries,
W.H. Freeman and Co., 1974.
AMS 1980 Subject Classification: 51 NXX
ANGULAR BOUNDARY VALUE, boundary value
along a non-tangential path - The value associated to a
complex function !(x) defined in the unit disc
201
ANGULAR BOUNDARY VALUE
D={ZEC: I Z I <1} at a boundary point ~=eill, equal
to the limit limf(z) =
zeS
rm
z--+\
of f(z) on the set of points of the angular domain
S(r,t:) = {z=re;<PED: I arg(eiO-z) 1<; -t:}
under the condition that this limit exists for all t:,
0<t:<'IT /2, and hence does not depend on t:. The term
is sometimes applied in a more general sense to func-
tions f(z) given in an arbitrary (including a higher-
dimensional) domain D; for S(~, t:) one takes the inter-
section with D of an angular (or conical) domain with
vertex ~EaD, with axis normal to the boundary at ~ aD
and with angle 'IT /2-t:, 0<t:<'IT /2.
References
[I] MARKUSHEVICH, A.I.: Theory offunctions of a complex variable,
1-2, Chelsea, 1977 (translated from the Russian).
[2) PIuwALOW, I.I. [1.1. PRIVALov): Randeigenschaften analytischer
Funktionen, Deutsch. Verlag Wissenschaft., 1956 (translated
from the Russian).
E.D. Solomentsev
Editorial comments. An angular boundary value is also
called a non-tangential boundary value. Cf. Boundary pro-
perties of analytic functions.
AMS 1980 Subject Classification: 30D40
ANHARMONIC RATIO - A synonym for double
ratio or cross-ratio.
AMS 1980 Subject Classification: 51 N15
ANISOTROPIC GROUP over a field k - A linear
algebraic group G defined over k and of k-rank zero,
i.e. not containig non-trivial k-split tori (cf. Splittable
group). Classical examples of anisotropic groups include
the orthogonal groups of quadratic forms that do not
vanish over k; and algebraic groups of elements of
reduced norm one in division algebras over k. If G is
semi-simple, and if the characteristic of k is zero, then
G is anisotropic over k if and only if Gk contains non-
trivial unipotent elements. (For the field of real
numbers or the field of p-adic numbers this is
equivalent to saying that Gk is compact.) The classifica-
tion of arbitrary semi-simple groups over the field k
reduces essentially to the classification of anisotropic
groups.
References
[IJ BOREL. A.: Linear algehraic r,roups, Benjamin. 1969.
[2J TITS, 1.: 'Classification of algebraic simple groups'. Pmc.
Symp. Pure Math .. Vol. 9. Amer. Math. Soc .. 1966. 33-62.
V. P. Platonov
AMS 1980 Subject Classification 20GXX .
ANISOTROPIC KERNEL - The subgroup D of a
semi-simple algebraic group G, defined over a fieJd k,
202
which is the commutator subgroup of the centralizer of
a maximal k-split torus SCG; D=[Zc(S),Zc(S)]. The
anisotropic kernel D is a semi-simple anisotropic group
defined over k; rank D = rank G - rankk G. The con-
cept of the anisotropic kernel plays an important role
in the study of the k-structure of G [I]. If D = G, i.e. if
rankk G = 0, then G is anisotropic over k; if D = (e), the
group G is called quasi-split over k.
References
[I] TITS, J.: 'Classification of algebraic simple groups', Proc.
Symp. Pure Math., Vol. 9, Arner. Math. Soc., 1966, 33-62.
[2] BOREL, A. and TITS, J.: 'Groupes roouctifs', Publ. Math. IHES
27 (1965),55-150.
VP. Platonov
AMS 1980 Subject Classification: 20GXX
ANNIHILATION OPERATORS - A family of closed
linear operators {a(j): fEH}, where H is some Hilbert
space, acting on a Fock space constructed from H (i.e.
on the symmetrization P(H) or anti-symmetrization
ra(H) of the space of tensors over H) such that on the
vector (j1 ... f,,)a Ef"(H), O'.=S, a, consisting of
the symmetrized (O'.=s) or anti-symmetrized (O'.=a) ten-
sor product of a sequence of elements fl' ... ,f" EH,
n = I, 2, ... , in H, they are given by the formulas:
a(j)(fl 0 ... 0 j,,)s = (1)
(f,;;)(j10 ... 0;;-1 0 ;;+1 0 ... 0f,,),.
in the symmetric case, and
a (j)(jl 0 ... 0 f"L = (2)
i=i
in the anti-symmetric case; the empty vector ~Era(H),
O'.=s, a (i.e. the unit vector in the subspace of constants
in ra(H is mapped to zero by a(j). In these formulas
(".) is the inner product in H. The operators
{a"(j): fEH} dual to the operators a(j) are called
creation operators; their action on the vectors
(jl'" f,,)a, O'.=s,a, n=1,2, ... , is given by the
formulas
a'(f)(f10'" 0f,,)a = (f0f I 0 0fn)a (3)
and
a'(j)fJ. =f
As a consequence of these definitions, for each n >0
the subspace r~(H);fn, O'.=S, a, the symmetrized or
anti-symmetrized n-th tensor power of H, is mapped by
a(j) into r~ _I and by a" (j) into r~ + 1(H).
In quantum physics, the Fock space ra(H), 0'.=.1', a,
is interpreted as the state space of a system consisting
of an arbitrary (finite) number of identical quantum
particles, the space H is the state space of a single par-
ticle. the subspace r~(H) corresponds to the states of
the system with n particles, i.e. states in which there are

just n particles. A state with n particles is mapped by
a(j) to a state with n - 1 particles ('annihilation' of a
particle), and by a * (j) to a state with n + 1 particles
(,creation' of a particle).
The operators a(j) and a * (j) form irreducible fami-
lies of operators satisfying the following permutation
relations: In the symmetric case (the commutation rela-
tions)
a(jl)a(!2)-a(j2)a(jl) = (4)
= a'(jI)a'(j2)-a'(j2)a'(jI) = 0,
and in the anti-symmetric case (the anti-commutation
relations)
a(jl)a(j2)+a(j2)a(jl) = (5)
= a'(jI)a'(j2)+a'(j2)a'(jI) = 0,
where E is the identity operator in P(H) or rU(H).
Besides the families of operators a(j) and a * (j), f E H,
described here, there exist in the case of an infinite-
dimensional space H also other irreducible representa-
tions of the commutation and anti-commutation rela-
tions (4) and (5), not equivalent to those given above.
Sometimes they are also called creation and annihila-
tion operators. In the case of a finite-dimensional space
H, all irreducible representations of the commutation
or anti-commutation relations are unitarily equivalent.
The operators {alf),a*(j):fEH} are in many con-
nections convenient 'generators' in the set of all linear
operators acting in the space ra(H), o:=s, a, and the
representation of such operators as the sum of arbitrary
creation and annihilation operators (the normal form of
an operator) is very useful in applications. The connec-
tion with this formalism bears the name method of
second quantization, cf. [1 J.
In the particular, but for applications important, case
in which H=L 2(R", d"x), 1'= 1,2, ... (or in a more
general case H=L 2 (M, Q), where (M, Q) is a measure
space), the family of operators
{a(j), a'lf): fEL2(R", d"x)} defines two operator-
valued generalized functions a(x) and a'(x) such that
a(j) = /a(x)f(x)d"x, a'(j) = /a'(x)f(x)dPx.
The introduction of a(x) and a * (x) turns out to be
convenient for the formalism of second quantization
(e.g. it allows one directly to consider operators of the
form
Xa(YI) ... a(Yn)d"xl ... d"xn d"yl ... d"Ym, n, m = 1, 2, ... ,
where K(xJ, ... ,xn;YJ, ... ,Ym) is a certain
'sufficiently-good' function), without having to recourse
to their decomposition as a series in the monomials
a'(jd' .. a'(j,,)a(gl) ... a(gm),
Al':Tl-COMMCTATIVE ALGEBRA
where
References
[I] BEREZIN, F.A.: The method of second quantization, Moscow,
1965 (in Russian).
[2] DOBRUSHIN, R.L. and MINLOS, R.A.: Uspekhi Mat. Nauk 32,
no. 2 (1977), 67-122.
[3] GARDING, L. and WIGHTMAN, A.: Proc. Nat. A cad. Sci. U.S.A.
40, no. 7 (1954), 617-626.
R.A. Minlos
Editorial comments.
References
[A1] GLIMM, J. and JAFFE, A.: Quantum physics, Springer, 1981.
[A2] BOGOLYUBov, N.N., LOGUNOV, A.A. and TODOROV, I.T.:
Introduction to axiomatic quantum field theory, Benjamin,
1975 (translated from the Russian).
[A3] BOER, J. DE: 'Construction operator formalism in many parti-
cle systems', In J. de Boer and G.E. Uhlenbeck (eds.): Stu-
dies in statistical mechanics, North-Holland, 1965.
AMS 1980 Subject Classification: 81 XX, 81 D05
ANNIHILATOR, left, of a set X in R - The set B,(X)
of all elementsy in R such thatyX=O. Here R is a ring
or a semi-group (or, generally, a groupoid) with a zero.
The right annihilator of a set X in R is defined in a
similar manner as the set
3rCx) = {zER: Xz=O}.
The set
3(X) = 3,(X) n ~:lr<X)
is the two-sided annihilator of X. In an associative ring
(or semi-group) R the left annihilator of an arbitrary
set X is a left ideal, and if X is a left ideal of R, then
8,(X) is a two-sided ideal of R; in the non-associative
case these statements are usually not true.
K.A. Zhevlakov
AMS 1980 Subject Classification: 20MXX, 13-XX, 16-
XX
ANNULAR DOMAIN - I) A doubly-connected
planar domain between two closed Jordan curves
without common points, one of each encloses the other.
2) For an annular domain with respect to a quadratic
differential, see Global structure of trajectories of a qua-
dratic differential.
E.D. Solomentsev
AMS 1980 Subject Classification: 54FXX, 53A04, 30-
XX
ANTI-COMMUTATIVE ALGEBRA - A linear algebra
over a field in which the identity
(*)
is valid. If the characteristic of the field differs from 2,
the identity (*) is equivalent with the identity
xy = - yx. All subalgebras of a free anti-commutative
algebra are free. The most important varieties of anti-
commutative algebras are Lie algebras, Mal'tsev alge-
203
ANTI -COMMUTATIVE ALGEBRA

bras and binary Lie algebras (cf. Lie algebra; Binary a ring A into a ring B that is an isomorphism of the
Lie algebra; MaI'tsevalgebra). additive group of A onto the additive group of B and
References for which (ab)<p=b<p'a<p (a, bEA).
O.A. Ivanova
[I] SHIRSHOV, A.I.: 'Subalgebras of free commutative and free
anti-commutative algebras', Mat. Sb. 34 (76), no. 1 (1954), 81- AMS 1980 Subject Classification: 13-01, 16AXX,
88 (in Russian). AT G 17-01
.. amov
W

AMS 1980 Subject Classification: 17-01, 15A69 ANTI-MOTION - A transformation of a pseudo-

Euclidean space which converts points at a real dis-
ANTI-CONFORMAL MAPPING, conformal mapping tance a from each other into points at a purely-
of the second kind - A continuous mapping w =f (z) of imaginary distance a. The motions and anti-motions of
a neighbourhood of a point Zo of the complex z-plane a pseudo-Euclidean space together form a group.
onto a neighbourhood of a point Wo of the complex w- A.B. Ivanov
plane which preserves the angles between the curves
AMS 1980 Subject Classification: 51 N25
passing through z 0 but changes the orientation. A func-
tion f (z) which produces an anti-conformal mapping is
ANTI-PARALLEL STRAIGHT LINES with respect to
an anti-holomorphic function. See also Conformal map-
two given lines m I and m2 - Two straight lines II and
ping.
12 which intersect m I and m 2 so that L I = L 2 (cf.
References Fig.).
[1] MARKUSHEVICH, A.I.: Theory offunctions of a complex variable,
1, Chelsea, 1977, Chapt. 2 (translated from the Russian).
E.D. Solomentsev
AMS 1980 Subject Classification: 30C35

ANTI-DISCRETE SPACE - A topological space in

which only the empty set and the entire space are open.
A.A. Mal'tsev
Editorial comments. Other frequently occurring names
for this topological space are indiscrete space and trivial
topological space.
AMS 1980 Subject Classification: 54AXX
ANTI-DISCRETE TOPOLOGY on a set - The topol-
ogy defined by the open base consisting of the empty
set and of the entire space.
A.A. Mal'tsev
AMS 1980 Subject Classification: 54AXX
ANTI-HOLOMORPHIC FUNCTION, anti-analytic
function - A function f (z) = u + iv of one or more com-
plex variables z = (z 1> ,zn) that is the complex con-
jugate of a holomorphic function f (z) = u - iv (cf. Ana-
lytic function).
E.D. Solomentsev
AMS 1980 Subject Classification: 30-01
ANTI-ISOMORPHISM OF PARTIALLY ORDERED
SETS - A bijective anti tone mapping of a partially
ordered set A into a partially ordered set B, i.e. a one-
to-one mapping <p: A --')B such that a <b (a, b EA)
implies a<p>b<p in B.
o.A. Ivanova
AMS 1980 Subject Classification: 06AXX
ANTI-ISOMORPHISM OF RINGS - A mapping <p of
If II and 12 are anti-parallel with respect to ml and
m2, then m I and m2 are also anti-parallel with respect
to II and 12 , In any quadrilateral inscribed in a circle,
any two opposite sides are anti-parallel with respect to
the other two sides. If the lines m I and m 2 intersect at
a point 0, one also says that II and 12 are anti-parallel
with respect to the angle m I Om 2. If the lines m I and
m 2 coincide, II and 12 are said to be anti-parallel with
respect to a straight line. A.B. Ivanov
AMS 1980 Subject Classification: 51 N05
ANTI-PARALLELOGRAM - A plane quadrangle in
which each pair of opposite sides is equal and anti-
parallel to the other two sides (cf. Anti-parallel straight
lines).
A.B. Ivanov
AMS 1980 Subject Classification: 51 M15, 51 M20
ANTI-PRISM - A semi-regular polyhedron (cf.
Semi-regular polyhedra) in which two parallel faces are
congruent regular n-gons, while the remaining 2n faces
are regular triangles (d. Fig.).
A.B. Ivanov
AMS 1980 Subject Classification: 51 M20

204

ANTI-SYMMETRIC TENSOR - See Tensor calculus.
AMS 1980 Subject Classification: 15A69, 53A45
ANTILOGARITHM of a number n - The number N,
denoted by ant log" n, the logarithm of which to the
base a is equal to n. Thus,
ant log" n = N = an,
or
log" N = n.
Antilogarithms are also called inverse logarithms.
AMS 1980 Subject Classification: 26A09
ANTINOMY, paradox - A situation in which two
mutual contradictory statements are demonstrated,
each one having been deduced by means that are con-
vincing from the point of view of the same theory.
As distinct from a sophism, which is a deliberately
presented erroneous conclusion containing a masked
error, an antinomy is usually an indication of deeper
deficiencies of the theory in question. The discovery of
an antinomy often leads to a thorough revision of the
theory as a whole, attracts attention to new effects, and
ultimately stimulates further studies. This feature of
antinomies attracted the interest of philosophers ever
since ancient times. One may recall, for example, the
important role played by antinomies in the philosophy
of E. Kant. A number of antinomies were studied in
antiquity under the name aporium. Here, two famous
antinomies, named after Zeno of Elea (5-th century
B.C.) will be quoted.
'Achilles and the tortoise'. This aporium deals with the
contradiction involved in certain attributes of motion,
and may be stated as follows. A runner (Achilles) is
standing at a point A, and a tortoise is standing at a
point B, 100 meters away from A. At the same moment
of time Achilles begins to run from A towards B, while
the tortoise begins to move away from B in the direc-
tion away from A at a speed which is, say, 100 times
less than that of Achilles. Experience teaches us that in
such a situation Achilles will overtake the tortoise after
a fairly short period of time. On the other hand, it is
also possible to conclude that Achilles will never over-
take the tortoise, and will not even reach the point B.
In fact, at the moment when Achilles has reached the
midpoint C I of the distance AB, the tortoise will have
moved away from B, though only by a short distance.
Next, Achilles will reach the midpoint C 2 of C I B, then
the midpoint C 3 of C 2B, etc. Throughout that time the
tortoise keeps moving away from B. In order to reach
B, Achilles must have been present at each point in the
infinite sequence of points C I, C 2 , .... However, it is
impossible to be present at an infinite number of points
in a finite period of time; therefore, Achilles will never
ANTINOMY
arrive at the point B and will never overtake the tor-
toise.
Paradoxes of this type are easily solved by methods
of the modem mathematical theory of motion. A
detailed analysis shows that such solutions materially
depend on the truth of the so-called Archimedean axiom
in the field of real numbers: For any real numbers
a, b >0 there exists a natural number n such that
an >b. Nevertheless, the problem illustrated by the
paradox is very real. It is possible to deny the conveni-
ence or the adequacy of the ordinary mathematical
model as a description of real motion. In order to
describe the concepts of infinitely small and infinitely
large physical magnitudes, it was repeatedly attempted
to construct a theory of real numbers in which the
Archimedean axiom would not be valid. In any case,
the theory of non-Archimedean ordered fields forms an
essential part of modem algebra. Such non-
Archimedean ordered fields - the non-standard real
straight line - play a determining role in non-standard
analysis.
The 'sandpile paradox' may be stated as follows.
Clearly, one grain of sand does not form a sandpile. If
n grains of sand still do not form a sandpile, it follows
that they will not form a pile after another grain of
sand has been added. Accordingly, no number of grains
of sand can form a sandpile.
This paradox can be solved in modern terms by say-
ing that the method of complete mathematical induc-
tion cannot be applied to undefined volumes such as a
'pile of sand grains'. In modern mathematics (second
half of the 20-th century) indefinite volumes are
employed in the foundations of mathematics to estab-
lish the consistency of classical theories, and the pro-
perties of such concepts are studied by exact methods.
It is possible to use methods of mathematical logic to
realize a situation in which mathematical induction in
its general form is not applicable to any natural series,
even though very many ordinary properties of natural
numbers remain valid in the process (the so-called
non-standard model of arithmetic).
However, it is the antinomies involving unusual for-
mulations of concepts which are undoubtedly of
greatest interest to mathematics. These are the so-called
logical and semantic antinomies (see also Skolem para-
dox).
Below three examples of logical antinomies are
presented.
Russell's paradox (B. Russel, 1902), also indepen-
dently discovered by E. Zermelo. Consider the follow-
ing property D of sets. Assume that a set X has the
property D if and only if it is not an element of itself.
Property D is valid for a very large majority of specific
sets used in mathematical considerations. E.g., neither
205
ANTINOMY
the set of natural numbers nor the set of all real
numbers are elements of themselves. Consider the set T
such that its elements are exactly those sets X that have
the property D. Now consider whether T E T or T f T is
true. If TET then, by definition, T has the property D,
i.e. T f T. Thus, one must postulate T f T. But this
again results in a contradiction, since if T f T, T has
the property D, so that TET.
The paradox can be avoided by postulating that all
that has been proved in this way is that there is no
such set T, i.e. that the property D does not character-
ize any set at all. However, this cannot be accepted,
since it is natural to suppose, from the point of view of
'naive' set theory, that any exactly-described property B
of objects defines a set C of objects which possess this
property. Russell's paradox is a strong blow to this
natural concept. One must conclude that certain very
simple properties must either be considered as impre-
cisely described, or else that there exist exactly-
described properties which do not define a set. This in
tum gives rise to several difficult questions. What are
'exactly' and what are 'inexactly' described properties?
Which properties define a set and which do not? May
not other properties which are extensively used in set
theory also lead to contradictions and have to be dis-
carded? Is it at least possible to reliably specify a
domain which is adequately protected from paradoxes,
and which is nevertheless sufficiently large to include
conventional mathematical practice?
With the construction of exact logico-mathematical
languages, the problem of the exact description of pro-
perties may be considered as being satisfactory solved.
However, the problem of description of criteria estab-
lishing the class of properties defining a set is still a
question that is far from solved. Even worse, modem
results in axiomatic set theory indicate that there is no
final solution to this problem. Russell's paradox made
a great impression on his contemporaries since it con-
cerned the very first stage in the study of set theory.
Nevertheless, there exist various ways of avoiding para-
doxes which, though hardly ultimate or natural, are
very convenient in practice and throw light both on the
nature of the paradox itself and on logical connections
of other principles of set theory. The axiomatic
approach to the fundamental principles of set theory
(cf. Axiomatic set theory) proved particularly success-
ful. Thus, the formal axiomatic system of
Zermelo - Fraenkel is at the time of writing (1970's) a
most useful axiomatic theory, which gives the most ade-
quate description of the 'non-paradoxical' part of the
practice of classical set theory.
The paradox of the 'village barber'. This is a variant
of Russell's paradox, formulated in the guise of an
everyday-life situation. (A slightly-different form of this
206
paradox is known as the Gonseth paradox.) It is the
case of a village barber who shaves those and only
those inhabitants of his village who do not shave them-
selves. Does the barber shave himself? A reasoning
similar to that given in the context of Russell's paradox
shows that he both shaves and does not shave himself.
This difficulty is easily overcome if one concludes that
there cannot be such a barber. The condition to be
satisfied by the barber is self-contradictory and there-
fore cannot be satisfied. Admittedly, the problem raises
the question of criteria of self-consistent properties but,
unlike in Russell's paradox, this is by no means a burn-
ing problem. It refers to a real-life situation, and most
such situations are never exactly formulated or reliably
defined. Moreover, self-consistency is not the only cri-
terion (nor probably even the most important one) of
acceptability of real-life reasoning. The situation is alto-
gether different in mathematical reasoning, which is
supposed to be much more precisely defined and more
reliable; self-consistency forms an important part of
such a reasoning.
Cantor's paradox (G. Cantor, 1899). Let M be the set
of all sets and let P (M) be the set of all its subsets. It
is clear from the definition of M that P (M) is included
in M. On the other hand, in accordance with a well-
known theorem of Cantor, the cardinality of P (M) is
larger than that of M, and for this reason P(M) cannot
be a subset of M. The conclusions which follow from
Cantor's paradox are roughly the same as those follow-
ing from Russell's paradox. In particular, it may be
thought that Cantor's paradox proves that the set M of
all sets does not exist. It is interesting to note in this
connection that there exist axiomatic systems of set
theory, such as the well-known 'New Foundations' sys-
tem of W. Quine, which establish the existence of M.
Cantor's paradox can be avoided in the 'New Founda-
tions' owing to the fact that Cantor's cardinality
theorem can be proved in this system only in a specific
form, which is insufficient to 'prove' the paradox. In
general, the concepts of set theory required to construct
Cantor's paradox are much more involved than those
needed to construct Russell's paradox (such as the con-
cept of a subset, the idea of a one-to-one correspon-
dence, etc.).
A paradox which resembles Cantor's paradox is the
so-called Burali - Forti paradox, which deals with the
class of all ordinal numbers. The order type of this
class must be larger than all the ordinals contained in
it.
Antinomies of a somewhat different type, the so-
called semantic antinomies, are also of major interest.
As distinct from logical antinomies, semantic
antinomies include terms such as 'truth', 'falsehood',
'denotes', 'defines' and others. Otherwise, the difference

between the two types of paradoxes, which was pro-
posed by F. Ramsey in 1926, is largely arbitrary.
Many semantic antinomies may be stated in a logical
form and vice versa. It was noted by Ramsey that
semantic antinomies cannot be constructed in ordinary
logico-mathematical theories for the simple reason that
such theories do not contain the semantic concepts
which are required in the formulation of semantic
antinomies. In this sense, semantic antinomies are
'safer' than logical antinomies. It should be noted, how-
ever, that there are modern theories containing con-
cepts for the formulation of certain semantic paradoxes,
in particular Richard's paradoxes (but such paradoxes
can always be avoided by using special devices). Below
two well-known semantic paradoxes are stated.
Richard's paradox in the formulation of G.D.W.
Berry (1906). Consider the set of natural numbers each
one of which can be uniquely defined by a meaningful
text containing not more than 1000 syllables. Obvi-
ously, the number of such texts is finite, since the col-
lection of all texts with 1000 syllables is finite. Consider
the smallest natural number that is not a member of
the set defined above.
The above paragraph is a meaningful text, containing
not more than 1000 syllables, which uniquely deter-
mines some natural number which, by definition, can-
not be characterized by a set of this type. Obviously,
this paradox can be avoided if this text is declared to
be not meaningful (or to be a text which does not
define a natural number), but in such a case, as before,
there arise difficult problems concerning the criteria for
a text to be meaningful, etc.
Eubulides' paradox (4-th century B.C.). A person
says: 'What I am saying now is a lie'. Is this statement
true or not? If it is true, its very meaning indicates that
it is false. If it is false, it immediately follows that it is
true. Many variants of this paradox are known - the
liar's paradox, Epimenides' paradox, etc. The idea of
this paradox forms the base of the proof of the well-
known COde! incompleteness theorem of formal
axiomatic theories.
In general, analysis of such paradoxes stimulated a
fundamental review of the views on the foundations of
mathematics and the development of many modern
ideas and methods of mathematical logic.
References
[1] FRAENKEL, A.A. and BAR-HILLEL, Y.: Foundations of set
theory, North-Holland, 1958.
[2] KLEENE. S.c.: Introduction to metamathematics, North-Holland,
1951.
A. G. Dragalin
AMS 1980 Subject Classification: OOA25, 03A05
ANTIPODES - Diametrically-opposite points on a
sphere. Borsuk's antipodal-point theorems apply [1]: 1)
APOLAR ]\;ETS
For any continuous mapping of the sphere sn into the
Euclidean space En there exist antipodes with a com-
mon imag~; 2) Any mapping of the sphere sn into
itself in which the images of antipodes are antipodes is
an essential mapping.
References
[II BORSUK, K.: 'Drei Siitze tiber die n -dimensionale euklidische
Sphare', Fund. Math. 20 (1933), 177-190. A V Ch k'-
.. ernavs II
Editorial comments. The first result mentioned above is
known as the Borsuk-Ulam theorem (on antipodes). The
following result also goes by the name of Borsuk's antipodal
theorem: There is no continuous mapping f of the (n + 1)-
ball B n +1 into the n-sphere sn such that f(x) = - f( - x), cf.
[A1],p.131.
References
[A 1] ISTRATEscu, V.l.: Fixed point theory, Reidel, 1981.
AMS 1980 Subject Classification: 55M20
ANTITONE MAPPING of partially ordered sets - A
mapping cf> of a partially ordered set A into a partially
ordered set B such that if a ~b (a, b EA), then acf>?3bcf>.
The dual concept to an anti tone mapping is an isotone
mapping.
o.A. Ivanova
AMS 1980 Subject Classification: 06A 10
APERIODIC AUTOMORPHISM of a measure space -
An automorphism T of a measure space such that its
periodic points, i.e. the points x for which Tk x =x for
some k >0, form a set of measure zero. The introduc-
tion of a special name for such transformations is due
to the fact that in certain theorems of ergodic theory
automorphisms with 'too many' periodic points are
considered as trivial exceptions (see [1 D.
References
[1] ROKHLIN, V.A.: 'Selected topics from the metric theory of
dynamical systems', A mer. Math. Soc. Transl. Series 2 49, 171-
240 (Uspekhi Mat. Nauk 4, no.2 (30) (1949), 57-128)
D. V Anosov
Editorial comments. The so-called Rokhlin- Halmos
lemma for periodic automorphisms is important tor the
approximation of an automorphism of a Lebesgue space by
periodic transformations (ct. Approximation by periodiC
transformations), see [A 1], p. 75, or [A2], p. 390.
References
[A 1] CORNFEL'D, LP., FOMIN, S.V. and SINAl, YA.G.: Ergodic
theory, Springer, 1982.
[A2] HALMOS, P.R.: Lectures on ergodic theory, Math. Soc. of
Japan, 1956.
AMS 1980 Subject Classification: 28-01, 28DXX
APOLAR NETS - Two nets given in the same
domain G of a two-dimensional manifold, where at
each point x E G the tangent directions of one net har-
monically subdivide the tangent directions of the other.
207
APOLAR NETS
Thus, for instance, the asymptotic net on a surface in a
Euclidean space is apolar with respect to the net of cur-
vature lines (cf. Curvature lines, net of).
VT Bazylev
AMS 1980 Subject Classification: 53A05, 53820
ApOLLONIUS PROBLEM - The problem of con-
structing, in a given plane, a circle tangent to three
given circles. The problem is solved by the method of
inversion. The circle representing the solution of this
problem is known as the Apollonius circle. The problem
was named after Apollonius of Perga (3-rd century
B.c.).
References
[I) Encyclopaedia of elementary mathematics, Vol. I, Moscow
Leningrad, 19511963 (in Russian).
A.B. Ivanov
Editorial comments.
References
[A 1) BAKER, H.F.: Principles of geometry, 4, F. Ungar, 1963.
AMS 1980 Subject Classification: 51 M15
APOLLONIUS THEOREM - I) The sum of the
squares of the lengths of conjugate semi-diameters of
an ellipse is constant and equal to the sum of the
squares of the lengths of its semi-axes.
2) The area of a parallelogram circumscribed around
an ellipse, with its sides oriented in opposite directions,
is constant and equal to the product of the lengths of
its diameters.
A.B. Ivanov
Editorial comments. Cf. [A 1] for a higher-dimensional
generalization.
References
[A1) BORSUK, K.: AnalytiC geometry, PWN, 1969.
AMS 1980 Subject Classification: 51 N1 0
ApOTHEM of a regular polygon - The segment (and
its length) of a perpendicular dropped from the centre
of the regular polygon onto any of its sides. The
apothem of a regular n-gon is equal to the radius of the
circle inscribed in it and is connected with the side of
the polygon, an' and with its surface area Sn by the
relations:
an = 2r" tan!!"', S" = nr;, tan!!....
n 11
The apothem of a regular pyramid is the height of its
(side) face.
AMS 1980 Subject Classification: 51 M15
ApPELL EQUATIONS - Ordinary differentia) equa-
tions which describe the motions of both holonomic
and non-holonomic systems, established by P.E. Appell
[1]. They are sometimes referred to as Gibbs - Appell
20~
equations, since they were first proposed by l.W. Gibbs
[3] for holonomic systems. The Appell equations in
independent Lagrange coordinates qs (s = 1, ... ,n)
have the form of second-order equations
~ - Q* i= I, ... ,k~n. (1)
3qi - i,
Here I N
S = 2' ~m.w;
,,=)
(m. and w. are the masses and the accelerations of the
N points of the system) is the energy of acceleration of
the system, which is so expressed that it contains the
second derivatives of the coordinates qi' i = 1, ... ,k,
only, the variations of which are considered as
independent; Q; are the generalized forces correspond-
ing to the coordinates q;, obtained as coefficients in
front of the independent variations oq; in the expres-
sion for the work of the given active forces F.
corresponding to virtual displacements or.:
N k
~F.~r. = ~Q; Oqi'
p:::: 1 i=1
In evaluating Sand Q; the dependent variables qj
(oqj) U= k + 1, ... ,n) are expressed in terms of the
independent velocities (variations) by solving the n - k
non-holonomic constraint equations (cf. Non-holonomic
systems), expressed in the generalized coordinates qs
(and by solving the equations for oqs obtained from
them). Differentiation with respect to the time t of the
expressions found for qj yields expressions for qj In
terms of qi'
Equations (1), together with the n - k equations of
the non-integrable constraints, form a system (of order
n + k) of n differential equations involving the n
unknowns qs'
For a holonomic system k = n, all velocities qs and
variations oqs are independent, Q; = Q;, and equations
(1) are a different notation for the Lagrange equations
(in mechanics) of the second kind.
Appell's equations in quasi-coordinates 1'" where
irr = ~ar,q"
" r=l, ... ,k, (2)
i=1
have the form
~=rr
air ro r = 1, ... , k ~11. (3)
Here S is the energy of acceleration, expressed in
terms of the second 'derivatives' 7Tr (with respect to the
time) of the quasi-coordinates, and I1r are the
generalized forces corresponding to the quasi-
coordinates. Equations (3), together with the n - k
equations of the non-integrable constraints and the k
equations (2), form a system of n + k differential equa-
tions of the first order with the same number of
 APPELL POLYNOMIALS

unknowns q" s=l, ... ,n, and ir" r=l, ... ,k. defined by "
oc A (z)
Appell's equations are the most general equations of A (t)e Zl = ~ _n_,_t n,
n=O n.
motion of mechanical systems.
References An(z) = nAn_\(z), n =],2, ... ,
[I] ApPELL,P.E.: 'Sur une forme generale des equations de la which, apart from normalization, are equivalent to
dynamique', C.R. A cad. Sci. Paris Ser. I Math. 129 (1899).
[2] ApPELL, P.E.: 'Sur une forme generale des equations de la those given above.
dynamique et sur Ie principe de Gauss', J. Reine Angew. Math. Appell polynomials of class A (I) are used to solve
122 (1900), 205-208. equations of the form:
[3] GIBBS, l.W.: 'On the fundamental formula of dynamics', A mer.
A(D)y(z) = f(z). (2)
1. Math. 2 (1879), 49-64.
vV Rumyantsev
The formal equality y(z)=f(z)/ A(D) for
Editorial comments.
OC Ck zk
References fez) = k~Ok!
[A1] WHIITAKER, E.T.: Analytical dynamics, Cambridge Univ.
Press, 1927, 258.
makes it possible to write the solution of (2) III the
AMS 1980 Subject Classification: 70GXX form

APPELL POLYNOMIALS - A class of polynomials

y(z) = ~ckA;(z),
k=O
over the field of complex numbers which contains where {AZ(z)} are Appell polynomials with the gen-
many classical polynomial systems. The Appell polyno- erating function e Zl / A (t). In this connection the
mials were introduced by P.E. Appell [l]. The series of expansion of analytic functions into Appell polynomi-
Appell polynomials {An(z)};'=o is defined by the for- als is of special interest. Appell polynomials also find
mal equality use in various problems connected with functional
equations, including differential equations other than
(I)
n=O (2), in interpolation problems, in approximation theory,
in summation methods, etc. (cf. [1] - [6]). For a more
where A (t)= 2:.'/:=oaktk is a formal power series with
general account of the theory of Appell polynomials of
complex coefficients ak> k =0, 1, ... , and ao==O. The
class A (I), and a number of applications, see [6].
Appell polynomials An(z) have an explicit expression in
The class A (1) con tains, as special cases, a large
terms of the numbers ak as follows:
number of classical sequences of polynomials. Exam-
ples, up to a normalization, are the Bernoulli polynomi-
als
te zt Bn(z)
The condition ao==O is tantamount to saying that the - t- = ~
x;
__ ,_tn;
e -] n =0 n.
degree of the polynomial An(z) is n.
There is another, equivalent, definition of Appell the Hermite polynomials
H (z)
polynomials. Let ezt-t /2
2
= ~
00
_n_ tn .
n ==0 n! '
the Laguerre polynomials
be a differential operator, generally of infinite order,
defined over the algebra P of complex polynomials in n=O
the variable Z = x + ry. Then etc. For many other examples, see [2] and [3], Vol. 3.
There exist various generalizations of Appell polyno-
A ( ) = A(D)zn n =0, I, ... ,
n Z n." mials, which are also known as systems of Appell polyno-
i.e. An(z) is the image of the function zn / n! under the mials. These include the Appell polynomials with gen-
mappingp =A (D)q, p, q EP. erating functions of the form
The class A (1) of Appell polynomials is defined as A(w)ezU(w) =
00

~Pn(z)wn,
the set of all possible systems of polynomials {An(z)} n=O
(3)
with generating functions of the form (1). To say that a
U(w) = ~Ckwk, c\*O,
system {Pn(z)} of polynomials (of degree n) belongs to k=\
the class A (1) amounts to saying that the relationships
and the Appell polynomials with the more general gen-
P~(z)=Pn-\(z), n=],2, ... ,
erating functions:
are valid. OC

The Appell polynomials of class A (I) are sometimes A(w)1/;(zU(w = ~qn(z)wn (4)
n=O

209
APPELL POLYNOMIALS
(see, for example, [2], [3], Vol. 3). If w(U) is the inverse
function to the function U(w), then the fact that the
system of polynomials {Pn(z)} ~=0 belongs to the class
of sequences of Appell polynomials with a generating
function of type (3) is equivalent to the validity of the
relations
w(D)Pn(Z)=Pn-I(Z), n=I,2, ... [D=~].
There are only five weighted orthogonal systems of
sequences of Appell polynomials on the real axis with
generating functions of the type (3); these include only
one orthogonal system with generating functions of the
type (I), which consists of Hermite polynomials with
the weight e - x' /2 on the real axis (cf. [7]).
For the expansion in series by Appell polynomials
with generating functions of the types (3) and (4), and
interconnections of these polynomials by various func-
tional equations see [2], [7], [8].
The class A (P), where p;;,.1 is an integer, of Appell
polynomials is defined as follows: It is the set of all
systems of polynomials {An(z)} for each of which the
(formal) representation
is valid. Here, wp = e 2'1Ti / P, and Ak(t), k = 0, ... ,p - I,
are formal power series, the free terms of which are
such that the degree of the polynomial An(z) is n. To
say that a sequence {Qn(z)} of polynomials of degree n
belongs to A (P) amounts to saying that the relations
DPQn(z) = QW)(z) = Qn -p(z), n =P,P + I, ... ,
are valid. For problems on the expansion of analytic
functions in series by Appell polynomials of class A (P),
see [9]. They are closely connected with the problem of
finding analytic solutions of functional equations of the
type p-I
~Ak(D)y(zw~) = fez).
k =0
Appell polynomials in two variables were introduced
by P. Appell [10]. They are defined by the equations:
I-y I-i
J m n (lX, Y, Y, , x, Y ) -- x y, (1- X)-
- ,)y+y
, (Y)nh)11
m, n =0. I, ... ,
III which it IS assumed that (Y)o = I,
(Y)n=y(y+I)" '(y+n-l) for n;;,.l; these Appell
polynomials are analogues of the Jacobi polynomials.
The Appell polynomials J m.n are orthogonal with the
weight
to any polynomials in two variables of degree lower
than m + n, over the domain T, where T is the triangle:
210
x >0, y>O, x + y < 1; however, they do not form a sys-
tem of orthogonal functions with the weight t (x, y) in
T (see, for example, [3], Vol. 2).
References
[I] ApPELL, P.E.: Ann. Sci. Ecole Norm. Sup. 9 (1880), 119-144.
[2] BOAS, R.P. and BUCK, R.C.: Polynomial expansions of analytiC
functions, Springer & Acad. Press (U.S.A. & Canada), 1958.
[3] BATEMAN, H. and ERDELYI, A.: Higher transcendental functions,
Vol. 1-3, McGraw-Hill, 1953-1955.
[4] WOOD, B.: 'Generalized Szitsz operators for the approximation
in the complex domain', SIAM J. Appl. Math. 17 (1969), 790-
801.
[5] SZEGO, G.: Orthogonal polynomials, Arner. Math. Soc., 1975.
[6] BOURBAKI, N.: Elements of mathematics. Functions of a real
variable, Addison-Wesley, 1976 (translated from the French).
[7] MEIXNER, J.: 'Orthogonale Polynomsysteme mit einer beson-
deren Gestalt der erzeugenden Funktion', J. London Math. Soc.
(1) 9 (1934),6-13.
[8] ANDERSON, CH.A.: 'Some properties of Appell-like polynomi-
als', J. Math. Anal. Appl. 19 (1967),475-491.
[9A] KAZ'MIN, YU.A.: 'Expansions in series of Appell polynomials',
Math. Notes 5, no. 5 (1969), 304-311. (Mat. Zametki 5, no. 5
(1969), 509-520)
[9B] KAZ'MIN, YU.A.: 'On Appell polynomials', Math. Notes 6, no.
2 (1969), 556-562. (Mat. Zametki 6, no. 2 (1969), 161-172)
[10] ApPELL, P.E.: Arch. Math. Phys. (1) 66 (1881),238-245.
Yu.A. Kaz'min
Editorial comments. For Appell polynomials in two vari-
ables an explicit biorthogonal system is known. There is
also an explicit system for the weight function (5), consist-
ing of products of two Jacobi polynomials and a power. See
[A1], p. 454, for references.
References
[A 1) KOORNWINDER, T.H.: 'Two-variable analogues of the classi-
cal orthogonal polynomials', in R. Askey (ed.): Theory and
applications of special functions, Acad. Press, 1975, pp.
435-495.
AMS 1980 Subject Classification: 33A70
APPELL TRANSFORMATION - A transformation of
the equations of plane motion of a material point under
the action of a force that varies with the coordinates
x, y of the point only. The Appell transformation is a
homographic transformation of the coordinates x, y into
the coordinates x I , Y 1 :
. ax+by+c a'x+b'y+c'
-a X XI = a"x+b"y+c YI = a"x+b"y+c'"
which involves the substitution of the time tI for t in
accordance with the formula
dt
k dtl = " " "2
(a x+b y+c )
.
As a result of the Appell transformation, the motion of
a point (x I, Y d, acted upon by a force F I (x 1, Y d that
depends only on the coordinates x l,y I, corresponds to
the motion of the point (x, y) acted upon by the force
(5 ) F(x, y). A transformation of the equations of motion of
a point (x, y) with the aid of formulas of general form
XI = cp(x.y). YI = l/;(x.y), dt l = ?I.(x.y)dt. (*)

yields the equations of motion of the point (x 1 ,y 1)
under the effect of a force that depends only on the
coordinates Xl, Y 1 only if (*) is a homographic
transformation. The Appell transformation is used to
solve Bertrand's problem on the determination of the
forces under which the trajectories of motion of
material points become conic sections. The Appell
transformation is named after P.E. Appell.
References
[1) ApPELL, P.E.: 'De I'homographie en mecanique', Amer. J.
Math. 12 (1890), 103-114.
(2) ApPELL, P.E.: 'Sur les lois de forces centrales faisant decrire a
leur point d'application une conique queUes que soient les con-
ditions initiales', Amer. J. Math. 13 (1891), 153-158.
L.N. Sretenskir
AMS 1980 Subject Classification: 34AXX
APPLICATE, z-AXIS - One of the Cartesian coordi-
nates of a point in the three-dimensional space.
AMS 1980 Subject Classification: 51 N20
APPROXIMATE COMPACfNESS - A property of a
set M in a metric space X requiring that for any x EX,
every minimizing sequence Yn EM (i.e. a sequence with
the property p(x, Yn)~P(x, M)) has a limit point Y EM.
Approximate compactness of a given set ensures the
existence of an element of best approximation for any
x EX. The concept of approximate compactness was
introduced [I] in connection with the study of Che-
byshev sets (d. Chebyshev set) in a Banach space,
which made it possible to describe convex Chebyshev
sets in certain spaces. In fact, let X be a uniformly-
convex smooth Banach space. For a Chebyshev set
M eX to be convex, it is necessary and sufficient that
it be approximately compact. It follows, in particular,
that the set of rational fractions with fixed degrees of
the numerator and the denominator is not a Chebyshev
set in the space Lp (1 <p < 00) if the degree of the
denominator is not smaller than one [I].
For subsequent studies on this subject see [2].
References
[Il EFIMOV, N.V. and STECHKIN, S.B.: 'Approximative compact-
ness and Cebysev sets', Soviet Math. Dok/. 2, no. 5 (1961),
522-524. (Dok/. Akad. Nauk SSSR 140, no. 3 (1961), 522-534)
[2) GARKAVI, A.L.: 'The theory of best approximation in normed
linear spaces', Itogi Nauk. Mat. Anal. 1967 (1969), 75-132 (in
Russian). v
.lu. N. SUbbotm
. X>Xo
AMS 1980 Subject Classification: 41 A65, 54D99
APPROXIMATE CONTINUITY - A generalization of
the concept of continuity in which the ordinary limit is
replaced by an approximate limit. A function f (x) is
called approximately continuous at a point Xo if
lim apf(x)
X---JoXo
= f(xo}.
APPROXIMATE DERIVATIVE
In the simplest case, f (x) is a real-valued function of
the points of an n-dimensional Euclidean space (in gen-
eral it is a vector-valued function). The following
theorems apply. 1) A real-valued function f (x) is
Lebesgue-measurable on a set E if and only if it is
approximately continuous almost-everywhere on E (the
Stepanov- Denjoy theorem). 2) For any bounded
Lebesgue-measurable function f (x) one has, at each
point xo,
~~ M(k) [f(x)d JL = f(xo),
where JL is the n-dimensional Lebesgue measure, R is an
n-dimensional non-degenerate segment containing xo,
and p is its diameter.
References
[1) SAKS, S.: Theory of the integral, Hafner, 1952 (translated from
the Polish).
G.P. To/stov
Editorial comments. For other references see Approxi-
mate limit.
AMS 1980 Subject Classification: 26A 15
APPROXIMATE DERIVATIVE - A generalization of
the concept of a derivative, where the ordinary limit is
replaced by an approximate limit. If, for a function
f (x) of a real variable x, the limit
. f(x)- f(xo)
lim ap "-'--'---=---
X-Xo x-xu
exists, then it is called the approximate derivative of the
functionf(x) at Xo and is denoted by f:p (xo). In the
simplest case f (x) is a real-valued function; in general
it is a vector-valued function. An approximate deriva-
tive may be finite or infinite. The classical rules for the
differentiation of a sum, a difference, a product, and a
quotient of functions apply to finite approximate
derivatives; the theorem on the differentiation of a
composite function does not apply in general. The con-
cept of an approximate derivative was introduced by
A.Ya. Khinchin in 1916.
Approximate Dini derivatives are defined by analogy
with ordinary Dini derivatives (d. Dini derivative):
Ad(XO) - lim sup from the right; Ad(XO) - lim inf
from the right; Ag(xo) - lim sup from the left; Ag(XO)
- lim inf from the left; for example,
A d (Xo ) -- lim f(x)- f(xo)
ap .
X-+Xo. X-Xo
The following Denjoy- Khinchin theorems apply. If a
real-valued function f (x) is finite and Lebesgue-
measurable on a set E, then, at almost any point of this
set, either f (x) has a finite approximate derivative, or
Ad(X) = Ag(x) = + 00, AAx) = Ag(X) = - 00.
If x
F(x) = jf(t)dt
a
211
APPROXIMATE DERIVATIVE
is a Denjoy- Khinchin integral, then F~p (x)= f(x)
almost-everywhere in the interval under consideration
(an ordinary derivative need not exist on a set of posi-
tive measure). This theorem explains the role played by
approximate derivatives in the theory of integrals.
There exist continuous functions without an ordinary
or an approximate derivative at any point of a given
interval.
Approximate partial derivatives of functions of
several real variables are considered as well.
References
[II SAKS, S.: Theory of the integral, Hafner, 1952 (translated from
the Polish).
G.P. Tolstov
Editorial comments. For other references see Approxi-
mate limit.
AMS 1980 Subject Classification: 26A24
APPROXIMATE DIFFERENTIABILITY - A general-
ization of the concept of differentiability obtained by
replacing the ordinary limit by an approximate limit. A
real-valued function f (x) of a real variable is called
approximately differentiable at a point x 0 if there exists
a number A such that
. f(x)-f(xo)-A(x-xo)
hm ap = O.
X--Ho X - Xo
The magnitude A (x - xo) is called the approximate dif-
ferential of f(x) at Xo. A function f(x) is approxi-
mately differentiable at a point x 0 if and only if it has
an approximate derivative lap (x 0) = A at this point.
Approximate differentiability of real functions of n real
variables is defined in a similar manner. For example,
for n = 2, f (x, y) is called approximately differentiable at
a point (xo,Yo) if
f(x,y)-f(xo,yo)-A(x-xo)-B(y-yo) - 0
lim ap - ,
x-xu'
y----I'O
P
where A and B are certain given numbers and
p= vex - xO)2 +(y _YO)2 . The expression
A (x - xo) + B(y -Yo) is called the approximate differen-
tial off(x,y) at (xo.Yo).
Stepanov's theorem: A real-valued measurable func-
tion f (x, y) on a set E is approximately differentiable
almost-everywhere on E if and only if it has finite
approximate partial derivatives with respect to x and to
y almost-everywhere on E; these partial derivatives
almost-everywhere on E coincide with the coefficients A
and B, respectively, of the approximate differential.
The concept of approximate differentiability can also
be extended to vector functions of one or more real
variables.
References
[I] SAKS, S.: Theory of the integral, Hafner. 1952 (translated from
the Polish).
G.P. To/stov
212
Editorial comments. For other references see Approxi-
mate limit.
AMS 1980 Subject Classification: 26A24
APPROXIMATE LIMIT - A limit of a function f (x)
as X-Ho over a set E for which Xo is a density point.
In the simplest case f (x) is a real-valued function of
the points of an n-dimensional Euclidean space; in the
more general case it is a vector function. The approxi-
mate limit is denoted by
lim apf(x).
X-Xo
In general, the existence of an ordinary limit does not
follow from the existence of an approximate limit. An
approximate limit displays the elementary properties of
limits - uniqueness, and theorems on the limit of a
sum, difference, product and quotient of two functions.
Let Xo be a density point of the domain of definition
of a real-valued function f (x). If the ordinary limit
limx - Ho f (x) exists, the approximate limit also exists
and is equal to it. The approximate upper limit of a
function f(x) at a point Xo is the lower bound of the
set of numbers y (including y = + (0) for which x 0 is a
point of dispersion of the set {x:f(xy}. Similarly,
the approximate lower limit of a function f (x) at a
point Xo is the upper bound of the set of points y
(including y = - (0) for which x 0 is a point of disper-
sion of the set {x: f (x )<y}. These approximate limits
are denoted, respectively, by
lim ap f(x) and lim ap f(x).
X-Xo X-Xo
An approximate limit exists if and only if the approxi-
mate upper and lower limits are equal; their common
value is equal to the approximate limit.
If x is real, one-sided (right and left) approximate
upper and lower limits are also used (xo must then be,
respectively, a right-hand or left-hand density point in
the domain of definition of the function). For the
approximate right upper limit the following notation is
used:
lim apf(x),
x-xu,
X>Xo
with corresponding notations for the other cases. If the
approximate right upper and lower limits coincide, one
obtains the right approximate limit; if the approximate
left upper and lower limits coincide, one obtains the left
approximate limit.
Approximate limits were first utilized by A. Denjoy
and A.Ya. Khinchin in the study of the differential
connections between an indefinite integral (in the sense
of Lebesgue and in the sense of Denjoy- Khinchin)
and the integrand (cf. Approximate continuity; Approxi-
mate derivative).

References
[I] SAKS, S.: Theory of the integral, Hafner, 1952 (translated from
the Polish). G P 7' I
. . 10
Editorial comments. A point of dispersion is defined as
follows: Let A be a set in W. Let 0A be the completely-
additive set function defined for measurable E by
0A(E) =1-1' (E n A), the outer measure of En A. Let x be
any pOint in Rn. The upper strong derivative and lower
strong derivative DoACx) and OOA(X) are called, respectively,
the upper outer density and lower outer density of A at x.
The point x is a point of density for a set A if the outer den-
sity of A at x is 1 and it is a point of dispersion if the outer
density of A at x is zero. If A is measurable, almost-all points
of A are pOints of density and almost-all points of its com-
plement are points of dispersion. The latter condition is also
sufficient for A to be measurable.
References
[A 1] BRUCKNER, A.M.: Differentiation of real functions, Springer,
1978.
[A2] MUNROE, M.E.: Introduction to measure and integration,
Addison-Wesley, 1953.
[A3] THOMSON, B.S.: Real functions, Springer, 1985.
AMS 1980 Subject Classification: 26A 15
APPROXIMATELY-COMPACf SET - A set having
the property of approximate compactness. A metric pro-
jection on any approximately-compact Chebyshev set is
continuous. Examples of approximately-compact sets
include boundedly-compact sets and closed convex sets
in the spaces Lp (0 <p < 00), and the set of rational
fractions in which the degrees of the numerator and
denominator are constant. In approximation theory and
in the theory of ill-posed problems frequent use is
made of spaces in which all closed sets are approxi-
mately compact.
References
[I] EFIMOV, N.V. and STECHKIN, S.B.: 'Approximative compact-
ness and Cebysev sets', Soviet Math. Dokl. 2, no. 5 (1961),
1226-1228. (Dokl. Akad. Nauk SSSR 140, no. 3 (1961), 522-
524)
[2] VLASOV, L.P.: 'Approximative properties of sets in normed
linear spaces', Russian Math. Surveys 28, no. 6 (1973), 1-66.
(Uspekhi Mat. Nauk 28, no. 6 (1973),3-66) Y N S bb
u. . u otm
AMS 1980 Subject Classification: 41 A65, 54099
APPROXIMATION - The replacement of mathemati-
cal objects by others which resemble them in certain
respects. It is possible to study in this way numerical
characteristics and qualitative properties of an object
by reducing the problem to the study of simpler or
more convenient objects (e.g. of objects with readily
calculable characteristics or with properties which are
known). Number theory deals with Diophantine
approximations, in particular with approximations of
irrational numbers by rational numbers. Approxima-
tions of curves, surfaces, spaces and mappings are stu-
APPROXIMA TIo,", BY PERIODIC TRA'-JSFORMATIO'-JS
died in geometry and in topology. Certain branches of
mathematics deal almost exclusively with approxima-
stov tions - e.g. the theory of approximation of functions
and the theory of numerical methods in analysis.
BSE-3
AMS 1980 Subject Classification: OOA25
APPROXIMATION BY PERIODIC TRANSFORMA-
TIONS - One of the methods of ergodic theory. Any
automorphism T of a Lebesgue space X with measure J.l
can be obtained as the limit of periodic automorphisms
Tn in the natural weak or uniform topology of the
space 2! of all automorphisms [1]. To characterize the
rate of approximation quantitatively one considers not
only the automorphisms Tn, but also finite measurable
decompositions of X which are invariant under to Tn, i.e.
decompositions of X into a finite number of non-
intersecting measurable sets Cn, I, . . . ,Cn,q,' which are
mapped into each other by Tn. The number
q,
d(T, Tn; ~n) = ~p.(TCn.i t1 TnC,,)
i=l
is an estimate of the proximity of Tn to T with respect to
~n; here A is the symmetric difference
At1B = (A \B) U (B \A).
If qn is given, it is possible to choose ~n and Tn (with
the above properties) such that d(T, Tn; ~n) is arbi-
trarily small [1]. The metric invariants of the automor-
phism T become apparent on considering infinite
sequences Tn and ~n such that for any measurable set A
there exists a sequence of sets An, each being the union
of some of the Cn,i, which approximates A in the sense
that
lim p.(A t1A n) =0
n~oo
('the decompositions ~n converge to a decomposition
into points'). If, in addition, d(T, Tn; ~n)<f(qn)' where
f (n) is a given monotone sequence tending to zero,
then one says that T admits an approximation of the
. first type by periodic transformations with rate f (n); if,
in addition, Tn permutes the sets Cn,i cyclically, then
one speaks of cyclic approximation by periodic transfor-
mations. For other variants see [2], [6], [7].
At a certain rate of approximation certain properties
of the periodic automorphisms Tn affect the properties
of the limit automorphism T. Thus, if T has a cyclic
approximation by periodic transformations with a rate
c / n, then, if c<4, T will be ergodic; if c<2, Twill
not be mixing; and if c < 1, the spectrum of the
corresponding unitary shift operator is simple. Certain
properties of T may be described in terms of the rate of
approximation. Thus, its entropy is equal to the lower
bound of the c's for which T admits an approximation
by periodic transformations of the first kind with a rate
213
APPROXIMATION BY PERIODIC TRANSFORMATIONS
of 2c /log2 n [2], [7]. Approximations by periodic
transformations were used in the study of a number of
simple examples [2], including smooth flows on two-
dimensional surfaces [8]. They served in the construc-
tion of a number of dynamical systems with unexpected
metric properties [2], [6], [7], or with an unexpected
combination of metric and differential properties [3],
[4].
The statement on the density of periodic automor-
phisms in 2.(, provided with the weak topology, may be
considerably strengthened: For any monotone sequence
j (n 0, the automorphisms which allow cyclic approx-
imations at the rate j(n) form a set of the second
category in 2.( [2]. Accordingly, approximations by
periodic transformations yield so-called category
theorems, which state that in 2.( (with the weak topol-
ogy) the automorphisms with a given property form a
set of the first or second category (e.g. ergodic sets are
of the second category, while mixing sets are of the first
category [I D.
Let X be a topological or smooth manifold, and let
the measure Jl be compatible with the topology or with
the differential structure. In the class of homeomor-
phisms or diffeomorphisms preserving Jl it is not the
weak topology, but other topologies that are natural.
Category theorems analogous to those for 2.( are valid
for homeomorphisms; for the history of the problem
and its present 'state-of-the-art' see [5].
References
[I) HALMOS, P.R.: Lectures on ergodiC theory, Math. Soc. of lapan,
1956.
[2) KAroK, A.B. and SrEPIN, A.M.: 'Approximations in ergodic
theory', Russian Math. Surveys 22, no. 5 (1967), 77-102.
(Uspekhi Mat. Nauk 22, no. 5 (1967), 81-106)
[3) ANosov, D.V. and KAroK, A.B.: 'New examples in smooth
ergodic theory. Ergodic diffeomorphisms', Trans. Moscow
Math. Soc. 23 (1970),3-36. (Trudy Moskov. Mat. Obshch. 23
(1970), 1-35)
[4) KAroK, A.B.: 'Ergodic perturbations of degenerate integrable
Hamiltonian systems', Math. USSR-Jzv. 7, no. 3 (1973), 535-
571. (Izv. Akad. Nauk SSSR Ser. Mat. 37 (1973),539-576)
(5) KAroK, A.B. and SrEPIN, A.M.: 'Metric properties of measure
preserving homeomorphisms', Russian Math. Surveys 25, no. 2
(1970), 191-220. (Uspekhi Mat. Nauk 25, no. 2 (1970), 193-
220)
[6) AKCOGLu, M.A., CHACON, R.V. and SCHWARTZBAUER, T.:
'Commuting transformations and mixing', Proc. A mer. Math.
Soc. 24, 637-642.
[7) SCHwARrzBAuER, T.: 'Entropy and approximation of measure
preserving transformations', Pacific J. Math. 43 (1972), 753-
764.
[8) KOCHERGIN, A.V.: 'On mixing in special flows over a shifting
of segments and in smooth flows on surfaces', Math. USSR-Sh.
25, no. 3 (1975),441-469. (Mat. Sh. 96, no. 3 (1975), 472-502)
D. V Anosov
Editorial comments. Contributions to the .foundation of
the theory of approximations were also made by VA Rokh-
lin (cf. [A2]).
If in an approximation by periodiC transformations one has
214
the following inequality for the sequences {~n}, {Tn}, where
Tn is periodic of order qn,
;=1
and Urn ~Ur in the strong topology for operators on
L 2 (X, p.), then one says that T admits an approximation of
the second type by periodic transformations with speed
fen). Reference [A1] is a basic and well-known one.
In this context one also speaks of partitions instead of
decompositions.
References
[A1) CORNFEL'D, I.P., FOMIN, S.V. and SINAT, YA.G.: Ergodic
theory, Springer, 1982, Chap!. 15; 16.
[A2) ROKHLIN, V.A.: 'Selected topics from the metric theory of
dynamical systems', Amer. Math. Soc. Trans!. Series 249,
171-240. (Uspekhi Mat. Nauk4, no. 2 (30) (1949),57-128)
AMS 1980 Subject Classification: 28DXX
APPROXIMATION IN THE MEAN - Approximation
of a given function jet), integrable on an interval [a, b],
by a function <l>(t), where the quantity
b
pIJ, </ = j I f(t)-cf>(t) I dt
a
is taken as the measure of approximation.
The more general case, when
b
pIJ,</ = jlf(t)-CP(t) Iqda(t) (q>O),
a
where aCt) is a non-decreasing function different from a
constant on [a, b], is called mean-power approximation
(with exponent q) with respect to the distribution da(t).
If a(t) is absolutely continuous and <l>(t)=a(t), then one
obtains mean-power approximation with weight <l>(t), and
if aCt) is a step function with jumps Ck at points tk in
[a, b], one has weighted mean-power approximation with
respect to the system of points {td with measure of
approximation
pIJ, </ = ~Ck I f(td -CP(td Iq.
k
These concepts are extended in a natural way to the
case of functions of several variables.
References
[I) GONCHAROV, V.L.: The theory of interpolation and approxima-
tion offunctions, Moscow, 1954 (in Russian).
(2) NIKOL'SKli, S.M.: Approximation offunctions of several variables
and imbedding theorems, Springer, 1975 (translated from the
Russian).
[3) RICE, l.R.: The approximation offunctions, I. Linear theory,
Addison-Wesley, 1964.
N.I. Korndchuk
VP. Motornyt
Editorial comments.
References
[A 1] TIMAN, A.F.: Theory of approximation of functions of a rea!
vanable, Pergamon, 1963 (translated from the Russian).
[A2] RIVLIN, TJ.: An introduction to the approXimation of func-
tions, Dover, reprint, 1981.
 APPROXIMATION OF A DIFFERENTIAL BOUNDARY VALUE PROBLEM BY DIFFERENCE BOUNDARY VALUE PROBLEMS
[A3] CHENEY, E.W.: Introduction to approximation theory, Chel-
sea, reprint, 1982.
AMS 1980 Subject Classification: 41 A25
APPROXIMATION OF A DIFFERENTIAL BOUN-
DARY VALUE PROBLEM BY DIFFERENCE BOUN-
DARY VALUE PROBLEMS - An approximation of a
differential equation and its boundary conditions by a
system of finite (usually algebraic) equations giving the
values of the unknown function on some grid, which is
subsequently made more exact by making the parame-
ter of the finite-difference problem (the step of the grid,
the mesh) tend to zero.
Consider the computation of a function U which
belongs to a linear normed space U of functions
defined in a given domain Du with boundary r, and
which is the solution of the differential boundary value
problem Lu = 0, Iu Ir = 0, where Lu = is a differential
equation, while Iu Ir =0 is the set of boundary condi-
tions. Let D hU be a grid (cf. Approximation of a dif-
ferential operator by difference operators) and let Uh be
the normed linear space of functions Uh defined on this
grid. Let [V]h be a table of values of the function v at
the points of D hU A norm is introduced into Uh so
that the equality
!iTo II [vlh II u, = II v II u
is valid for any function v E U. The problem of comput-
ing the solution U is replaced by a certain problem
!l'huh =0 for the approximate computation of the table
[u lh of values of U at the points of D hU Here,!l'hUh is
a certain set of (non-differential) equations for the
values of the grid function Uh E Uh
Let Vh be an arbitrary function of Uh , let !l'hvh =CPh,
and let cI>h be the normed linear space to which CPh
belongs for any Vh E Uh One says that the problem
!l'hUh is a finite-difference approximation of order p of
the differential boundary value problem Lu =0,
Iu Ir =0, on the space of solutions U of the latter if
II.Ph [u1, 11<1>, = O(hP).
The actual construction of the system !l'huh involves a
separate construction of its two subsystems LhUh =0
and IhUh Ir, =0. For LhUh =0 one uses the difference
approximations of a differential equation (cf. Approxi-
mation of a differential equation by difference equa-
tions). The complementary equations IhUh Ir, =0 are
constructed using the boundary conditions 1Uh Ir =0.
An approximation such as has just been described
never ensures [2] that the solution Uh of the finite-
difference problem converges to the exact solution u,
i.e. that the equality
lim II [ulh -Uh II u,
h-+O
= difference analogues started in the works [1], [2] and
is valid, no matter how the norms in Uh and cI>h have
been chosen.
The additional condition, the fulfilment of which in
fact ensures convergence, is stability [3], [5] - [8], which
must be displayed by the finite-difference problem
!l'hUh =O. This problem is called stable if there exist
numbers ~>o and O<ho such that the equation
!l'hzh =CPh has a unique solution Zh E Uh for any CPh EcI>h,
II cf>h II <~, h <h o, and if this solution satisfies the ine-
quality
II Zh - Uh II u, ..; C II <Ph II <1>"
where C is a constant not depending on h or on the
perturbation CPh of the right-hand side, while Uh is a
solution of the unperturbed problem !l'huh =0. If a
solution U of the differential problem exists, while the
finite-difference problem !l'hUh approximates the dif-
ferential problem on solutions U of order p and is
stable, then one has convergence of the same order, i.e.
II [U1,- Uh Ilu, = O(hP).
For instance, the problem
L(u) - ~~ - ~~ = 0, t>O, - OO<X<OO,}
(1)
IU/ r = u(O,x)-l/;{x) = 0, -oo<x<oo,
where #.x) is a given function with a bounded second-
order derivative, can be approximated, for a natural
definition of the norms, by the finite-difference prob-
lem
= E cI>h,
where ui!z is the value of Uh at (xm' In) = (mh, m), r=rh,
r =const. If the norm of CPn is taken to be the upper
bound of the moduli of the right-hand sides of the
equations which constitute the system !l'h Vh =CPh,
Vh E Uh , then the approximation of problem (1) by
problem (2) on solutions U is of the first order. If r> 1,
there is no convergence, whatever the norm. If r';;;;; 1
and the norm is
II Uh II U. = sup
m.n
I u;:' I,
the problem is stable, so that there is convergence [2],
[3]:
II [u1, -Uh II U. = O(h).
The replacement of differential problems by differ-
ence problems is one of the most universal methods for
the approximate computation of solutions of differen-
tial boundary value problems on a computer [7].
The replacement of differential problems by their
[4], and is sometimes employed to prove that the dif-
215
APPROXIMATION OF A DIFFERENTIAL BOUNDARY VALUE PROBLEM BY DIFFERENCE BOUNDARY VALUE PROBLEMS

ferential problem is in fact solvable. This is done as fol- For instance, the equation
lows. It is proved that the set of solutions Uh of the _ d 2u du
difference analogue of the differential boundary value Lu = dx2 +p(x) dx +q(x)u =0
problem is compact with respect to h, after which a is approximated with second-order accuracy by the
proof is given that a solution U of the differential boun- finite-difference equation
dary value problem is the limit of a subsequence Uh,
um+I-2um+Um-1
which converges as hk~O. If this solution is known to LhUh h2 +
be unique, then not only the subsequence, but also the
Um+I-Um-1
entire set of Uh converges to the solution U as h~O. +p(xrn) 2h +q(xm)um = 0,
References
[I] LYUSTERNIK, L.A.: 'Dirichlet's problem', Uspekhi Mat. Nauk 8
where the grids DhU and DhF consist of points Xm = mh,
(1940), 125-124 (in Russian). m is an integer, and Um is the value of the function Uh
[2] COURANT, R., FRIEDRICHS, K. and LEWY, H.: 'Ueber die par- at X m Again, the equation
tiellen Differenzengleichungen der mathematischen Physik',
Math. Ann. 100 (1928). L =
U-
~-
at
a2 u = 0'
ax2
[3] GODUNOV, S.K. and RYABEN'KIl, V.S.: The theory of difference
schemes, North-Holland, 1964.
[4] PETROVSKIl, I.G.: 'New existence proofs for the solution of the is approximated by two different difference approxima-
Dirichlet problem by the method of finite differences', Uspekhi tions on smooth solutions:
Mat. Nauk 8 (1940), 161-170 (in Russian).
Unn+1 -Umn Unm+1 -
2n+n
Um Urn -I
[5] RYABEN'KIl, V.S.: 'On the application of the method of finite
differences to the solution of the Cauchy problem', Dokl. Akad.
L hIl) Uh =
-

T
- h2
=0
Nauk SSSR 86, no. 6 (1952), 1071-1073 (in Russian). (explicit scheme), and
[6] RJABENKI, V.S. [V.S. RYABEN'KIi] and FILlPOV, A.F. [A.F.
Un + 1 - un un + 1 _ 2un + 1 + un + 1
FILlPPOV]: Uber die stabilitiit von DifJerenzgleichungen,
Deutsch. Verlag Wissenschaft., 1960 (translated from the Rus-
L~2)Uh m rn _ m + 1 rn m -I =0
T h2
sian).
[7] SAMARSKIl, A.A.: Theorie der DifJerenzverfahren, Akad. VerI. (implicit scheme), in which the grids DhU and DhF con-
Geest u. Portig K.-D., 1984 (translated from the Russian). sist of the points (xm, tn)=(mh, n7), T=rh2, r=const,
[8] FILlPPOV, A.F.: 'On stability of difference equations', Dokl.
Akad. Nauk SSSR 100, no. 6 (1955), 1045-1048 (in Russian).
m and n are integers, and U:;, is the value of the func-
VS. Ryaben'kit tion Uh at the point (xm, tn) of the grid. There exist
finite-difference operators Lh which represent an
Editorial comments. For additional references, see the approximation to the differential operator L especially
additional references to Approximation of a differential well only on solutions U of the equation Lu = 0 and less
operator by difference operators. well on other functions. For instance, the operator L h ,
AMS 1980 Subject Classification: 65LXX, 65MXX
_ Urn + 1 - Urn [ h Urn +U] [ h]
Lhun = h -</> xrn + 2 ' - 2 - = Ji, xm+2'
APPROXIMATION OF A DIFFERENTIAL EQUA-
TION BY DIFFERENCE EQUATIONS - An approxi- where U=Um +hcj>(xm ), is a first-order approximation
mation of a differential equation by a system of alge- (with respect to h) of the operator L,
braic equations for the values of the unknown functions Lu ~~ -</>Cx, u) = j(x),
on some grid, which is made more exact by making the
parameter (mesh, step) of the grid tend to zero. on arbitrary smooth functions u(x) and it is a second-
Let L, Lu = f, be some differential operator and let order approximation on solutions of the equation
L h, LhUh = Ji" Uh E Uh, Ji, EFh, be some finite-difference Lu = 0 (it is assumed that the function u is sufficiently
operator (cf. Approximation of a differential operator by smooth). When finding numerical solutions of boun-
difference operators). One says that the finite-difference dary value problems for the differential equation Lu = 0
expression LhUh =0, OEFh, is an approximation of with the aid of the finite-difference equation LhUh =0, it
order p to the differential equation Lu =0 on solutions is the approximation property of Lh when applied to
U if the operator Lh approximates the operator L on solutions U of Lu =0 that is significant and not its
solutions U of order p, i.e. if approximation property when applied to arbitrary

II f, = O(hP).
I Lh[u lh
The simplest example of the construction of a finite-
difference equation Lh Un = 0 which approximates the
differential equation Lu = 0 on solutions u, consists in
replacing each derivative in the expression Lu by its
finite-difference analogue.
smooth functions. For a large class of differential equa-
tions and systems of equations there exist methods for
constructing difference equations approximating them,
which also meet various additional conditions: stability
of the solution Uh with respect to rounding-off errors
such as are permitted in the calculations; the validity of
certain integral relations for Uh which hold for the solu-

216
 APPROXIYlATIO:--; OF A DIFFERENTIAL OPERATOR BY DIFFERENCE OPERATORS

tion u of the differential equation; the permissibility of function Lu from the table of values [u 1r of the func-
using arbitrary grids DhU and DhF (this is important in tion u and for the approximation of a differential equa-
calculating the motion of a continuous medium); a lim- tion by difference equations.
itation on the number of arithmetical operations There are two principal methods for constructing
required to compute the solutions; etc. operators Lh approximating L.
The approximation of a differential equation by In the first method Lh[u 1h is defined as the result of
difference equations is an element of the approximation applying the differential operator L to a function in U,
of a differential boundary value problem by difference obtained by some interpolation formula from the grid
boundary value problems in order to approximately cal- function [u 1h'
culate a solution of the former. The second method is as follows. In the domain DF
References of definition of a function fin F one introduces a grid
[I) GoDUNOV. S.K. and RYABEN'KIi, V.S.: The theory of difference DhF , and considers the linear space Fh of grid functions
schemes, North-Holland, 1964 (in Russian). defined on DhF . The operator Lh[u 1r is constructed as
[2) SAMARsKIi, AA: Theorie der Differenzverfahren, Akad. VerI.
Geest u. Portig K.-D., 1984 (translated from the Russian). the product of two operators: an operator which con-
[3) GoDUNOV, S.K., ET AL.: Numerical solution of multi-dimensional verts the function [u 1r into the grid function fh in Fh ,
problems of gas dynamics, Moscow, 1976 (in Russian). i.e. into a table of approximate values of f, and an
[4) SAMARSKIi, AA and PoPov, Yu.P.: Difference schemes of gas
dynamics, Moscow, 1975 (in Russian).
operator which extends fh from DhF to the entire
VS. Ryaben'kif domain Df . For instance, in order to approximate the
Editorial comments. For additional references, see the differential operator
additional references to Approximation of a differential
L -- ~
dx' du -- ji()
dx x, 0,;:::: ,;:::: 1
""'Z""',
operator by difference operators.
AMS 1980 Subject Classification: 65L99, 34A99 one constructs the grid D hU consisting of points Xb
k=O, ... ,N,
APPROXIMATION OF A DIFFERENTIAL OPERA-
TOR BY DIFFERENCE OPERATORS - An approxima-
= Xo < ... < Xk < Xk+l < ... < XN = 1,
tion of the differential operator by parameter- max(Xk+l -Xk)
k
= h,
dependent operators such that the result of their appli-
and a grid DhF consisting of the points
cation to a function is determined by the values of this
function on some discrete set of points - a grid - x; =Xk+8(Xk+l-Xk), k=0, ... ,N-1,
which become more exact as its parameter (mesh, step .; ; 8 .;;;; 1, 8 = canst
of the grid) tends to zero. The values of the operator Lhlu 1h at the points x; are
Let L, Lu = f, be a differential operator which con- defined by the equations:
verts any function u of a class of functions U into a
- fi( ') -
function f of a linear normed space F. Let Du be the L hUh
[ 11 U(Xk+l)-U(Xk)
. - h Xk - ,
x=x, Xk+l- Xk
domain of definition of the functions in U, and let
there be some discrete subset in Du - a grid DhU - k=O, ... ,N-I.
which 'becomes more dense' as h~O. Consider the set Thereafter, the definition of Lh[u 1h is piecewise linearly
Uh of all functions [u 1r defined on the grid only and extended outside DhF with possible breaks at the points
coinciding with u in the points of the grid. A difference Xk, k= 1, ... ,N -2, only.
operator is defined as any operator Lh that converts the Let the norm in F be defined by the formula
grid functions in Uh into functions fh in F. One says
that the operator L h, Lhlu 1h = fh, represents an order p
II <I> II F = sup
x
I </>(x) I
approximation to the differential operator L on U if for Then, on the class of functions U with a bounded third
any function u E U derivative, for (}=O and (}=h /2 the operator Lh
I Lu - Lh[u1h IIF ~ 0, represents an order 1, respectively 2, approximation to
L =d / dx. On the class U of functions with bounded
I Lu-Lh[u1. IIF .;;;; ch P , c=c(u)=const, second derivatives, the representation is of order 1
as h~O. Occasionally, an approximation is understood only, for any (}E[O, 11.
to be the equality The task of approximating a differential operator by
finite-difference operators is sometimes conditionally
considered as solved if a method is found for the con-
in the sense of some weak convergence. The approxi- struction of the grid function
mation of a differential operator by difference opera-
tors is used for an approximate computation of the Lh[U1.1 = Ji, E Fh,
D"

217
APPROXIMATION OF A DIFFERENTIAL OPERATOR BY DIFFERENCE OPERATORS
determined at the points of DhF only, while the task of
completing the function fh everywhere on DF is
ignored. In such a case the approximation is defined by
considering the space Fh as normed, and by assuming,
for the grid and for the norm, that for any function
fEF, the function {f}h EFh , which coincides with f at
the points of DhF , satisfies the equation
~ II {fh IIF. = IIfIiF'
The operator Lh is understood to be an operator from
Uh in Fh , and one says that Lh represents an order p
approximation to L on U if, for h~O,
II {Lu h- Lh[u1.IIF. ~ 0,
I {Luh - Lh[U1. IIF. ~ ch P
In order to construct an operator Lh which is an
approximation to L of given order on sufficiently
smooth functions, one often replaces each derivative
contained in the expression L by its finite-difference
approximation, basing oneself on the following fact.
For any integers i, j and for any ko, 2ko + l-;;;:;'i + j, in
the equation
ko
h-j L cku(x+kh) = uV)(x)+l(x,h,Ck o" "
k= -k o
it is possible, by using the method of undetermined
coefficients and Taylor's formula, to select numbers Ck
not depending on h, so that for any function u (x) with
j + r (r ~ i) bounded derivatives, an inequality of the
type
I l(x,h,Ck o" " ,Ck o ) I ~ Aijsup
t
I uv+r)(t) Ih r ,
where Ai) depends only on i and j, is valid. As an
example, suppose one constructs an approximating
operator for the Laplace operator Ll,
!::"u -
a2 u
-2 +-2
a2 u = f(x,y),
ax ay
if Du is the closed square I x I ~ 1, Iy I ~ 1, and DF is
its interior I x 1<1, Iy 1< l. Assume that h = 1/ N
where N is a natural number, and construct the grid
with points
(x,y) = (mh, nh), I mh I ~ 1,
which belong to DhU ' The points
(x,y) = (mh,nh), I mh I ~ 1,
then belong to D hF , for integers m and n. Since
1 h
J;2[r(X+ )-2y(x)+y(x-h)] = y "{x)+U?
Ll can be approximated with second-order accuracy on
a space of sufficiently smooth functions by the finite-
difference operator Lh if one puts, at the points of
218
+ u""m + 1 - 2um,/I + Um,n - 1 f,
h2 = mn,
where um,n and fm,n are the values of the functions [u lit
and fh at the point (mh, nh).
There are also other methods of constructing opera-
tors Lh which are approximations to the operator L on
the space of solutions u of the differential equation
Lu = 0, and which satisfy additional conditions.
References
[I) FILIPPOV, A.F.: 'On stability of difference equations', Dokl.
Akad. Nauk SSSR 100, no. 6 (1955), 1045-\048 (in Russian).
(2) BEREZIN, I.S. and ZHIDKOV, N.P.: Computing methods, Per-
gamon Press, 1973 (translated from the Russian).
V.S. Ryaben'kit
Editorial comments. The approximation of a differential
operator by difference operators is an ingredient for both
the approximation of a differential equation by difference
equations and for the approximation of a differential
boundary value problem by difference boundary value
problems, and is therefore extensively treated in the litera-
ture on (finite-) difference methods for ordinary and partial
differential equations. The references listed below not only
provide discretizations of differential equations and boun-
dary value problems, but also the solution of these prob-
,Ck o )' lems. References [A1] - [A3] , [A5], [A6] are introducing
textbooks, whereas [M], [Al], [AB] and [A9] also present
more advanced material.
References
[A1] AMEs, W.F.: Numerical methods for partial differential equa-
tions, Nelson, London, 1969.
[A2] FORSYTHE, G.E. and WASOW, W.R.: Finite difference methods
for partial differential equations, Wiley, 1960.
[A3] GARABEDIAN, P.R.: Partial differential equations, Wiley, 1964.
[A4] GODUNOV, S.K. and RYABEN'KIT, V.S.: The theory of differ-
ence schemes, North-Holland, 1964 (translated from the
Russian).
[A5] LAMBERT, J.D.: Computational methods in ordinary differen-
tial equations, Wiley, 1973.
[A6] MITCHELL, A.R. and GRIFFITHS, D.F.: The finite difference
method in partial differential equations, Wiley, 1980.
[A7] RrCHTMEYER, R.D. and MORTON, K.W.: Difference methods
for initial value problems, Wiley, 1967.
[A8] SAMARSKIT, A.A.: Theorie der Differenzverfahren, Akad. VerI.
Geest u. Portig K.-D., 1984 (translated from the Russian).
[A9] STETIER, H.J.: Analysis of discretization methods for ordinary
differential equations, Springer, 1973.
AMS 1980 Subject Classification: 65D25, 65LXX,
I nh I ~ 1, 65MXX
APPROXIMATION OF FUNCfIONS - The replace-
I nh I ~ 1, ment, according to a definite rule, of a function by a
function <p close to it in some sense and belonging to a
set 9C (the approximating set) that is prescribed in
h 2 (4)
(~), advance. It is assumed that f is defined on a set Q in
m-dimensional Euclidean space (the real axis being a
particular case) in which the approximation is carried
out; f can either be given explicity in terms of elemen-
tary functions or f can be the solution of some equa-
tion. If only incomplete information about f is avail-
able, the approximation problem is in fact a matter of

approximating the whole class of functions defined by
the information available.
In practice, the necessity to approximate functions
arises in various situations when it is needed to replace
a function by a smoother function or by one which is
simpler and more convenient in computations, or when
it is required to establish a functional dependence on
the basis of experimental data, etc.
In the general problem of the approximation of func-
tions it is usually possible to single out the following
more particular problems: the choice of the approxi-
mating set 9C; the choice of the measure of the error of
approximation; the choice of the method of approxima-
tion, e.g. the rule according to which a function cP in 9C
is made to correspond to f; and the investigation and
estimation of the error of the approximation.
With respect to the choice of the approximating set 9C,
apart from the necessity of ensuring the required accu-
racy of the approximation, one is guided by the aim to
deal with structurally simple functions cp that are con-
venient in computations and that can be submitted to a
priori specified conditions, e.g. related to smoothness.
The classical tools for approximation consist of alge-
braic polynomials (if Q is a bounded closed set) and
trigonometric polynomials (in the periodic case) in one
or more variables. Their widespread application as
approximating sets is due, in particular, to the fact that
it is in principle possible to approximate any continu-
ous function by means of algebraic of trigonometric
polynomials with any degree of accuracy specified in
advance. The accuracy of the approximation can be
increased at the expense of increasing the degree of the
polynomial; however, this complicates the approximat-
ing process and increases the difficulties of applying it.
In practice one takes subspaces of algebraic or tri-
gonometric polynomials of a fixed degree as approxi-
mating sets and one aims at obtaining the accuracy
required while keeping the degrees of the polynomials
used as low as possible. A more general, and at the
same time more flexible, approximation tool is obtained
by considering generalized polynomials
!f>(t) = Clf[>l(t)+ ... +CNf[>N(t),
where {CPI, ... ,CPN} is a system of linearly independent
functions which should be chosen such that the condi-
tions of the problem at hand and the a priori condi-
tions on cP are taking into account.
In many problems splines (cf. Spline) have been
found to be more natural and more convenient from
the viewpoint of computation than classical polynomi-
als. If
fjw = {a =7'0<7'1 < ... <7'N=b}, N~2,
is a fixed subdivision of the interval [a, b), then a (poly-
nomial) spline of degree r and defect k (k = 1, . . . ,r)
APPROXIMA nON OF FUNCTIONS
with respect to tlN is a function s composed of alge-
braic polynomials of degree r 'tied together' at the
points TI, . . . ,TN -I in such a way that on the whole
interval [a, b) the spline s is continuous together with
its derivatives up to and including order r - k. Thus
SEC -k[a, b) and on each interval (Tj -I> Tj)
(i = 1, ... ,N) s is an algebraic polynomial of degree at
most r. For instance, a polygonal line with corners at
the nodes (knots, joints) Tj is a spline of degree one and
defect one; a function s that is continuously differenti-
able on [a, b) and that coincides on [Tj - I> Til,
i = I, ... ,N, with a cubic polynomial is a cubic spline
of defect 2, etc. Splines in two or more variables are
similarly defined. Having finite smoothness properties,
splines have a greater local flexibility than polynomials;
a change in the values of a spline in an interval (a, f3)
does not much affect (if at all) its behaviour outside
(a, f3). Apart from the simplicity of their implementa-
tion in computers, the advantages of splines show up,
in particular, when the information about the function
f to be approximated has a discrete character, for
instance, when the values of f and, possibly, of some of
its derivatives, are given at some points.
If f has singularities or if the approximation is car-
ried out on an unbounded domain, a convenient tool
for approximation consists of rational fractions p / q,
where p and q are algebraic polynomials (rational
approximation, Pade approximation). Non-periodic
functions defined on the whole real axis can also be
approximated by entire functions of exponential type.
The measure of the error of approximation, pIJ, cp), is
usually chosen by taking into account the condition of
the concrete problem and the information available
about the function to be approximated. Most fre-
quently it is a matter of choosing a function space con-
taining f in whose metric it is expedient to estimate the
error of approximation. If
p.(j, f[ = sup I f(t)-!f>(t) I,
teQ
then one deals with uniform, or Chebyshev, approxima-
tion, while if
p.(j,f[ = jlf(t)-!f>(t) IPdt, p>O,
Q
then one has mean-power approximation, or 4-
approximation, which in the case of p = I is called
approximation in the mean. A particular important case
is that of p =2, mean-square approximation: the error of
the best approximation of the function f by a finite-
dimensional subspace can be expressed exactly in terms
of certain determinants. In some problems the
requirements of proximity (closeness) of the functions f
(1)
and cP are different at different points; in order to take
into account this lack of homogeneity a weight function
p(t);;a.O is introduced and weighted approximation is
219
APPROXIMATION OF FUNCTIONS

considered, where the measure of the error of approxi- If W is a linear manifold and if it contains a system of
mation is given by functions cp], ... ,CPN such that cp;(1;) = 1 and CP;(tk)=O
pif, </ = ~~8 I f(t)-<J>(t) I p(t) when ki=i, then the function
N
or
pif, </ = f I f(t)-<J>(t) IPp(t)dt, p>O.
</>(f, t) = ~f(tk}</>k(t)
k=l
Q
The weight function allows one also to see to it that the belongs to ~ and satisfies (2); it yields an interpolation
error is finite when, for instance, f is unbounded. If the method of approximation and is obviously linear. It may
error of approximation has to take into account the be required to choose the function <P(f, t) in such a way
closeness of f and </> only at some individual points tko that at the points tk it coincides not only with f, but
k = 1, ... ,N, in Q, then as JL(f, </ one may take one of also with some its derivatives; this case is called inter-
the quantities polation with multiple nodes. If Q = [a, b) and
a.,;;;;;tl < ... <tN.,;;;;;b, then there is a unique algebraic
polynomial of degree N - 1, and in the periodic case
or N
~ If(td-<J>(tk) IP, p>O,
(b-a=27T, tN =t2m-l<b) - a unique trigonometric
k=l polynomial of degree n -1, coinciding with f at the
into which weight coefficients can also be introduced. points tk' Multiple interpolation is achieved with Her-
When deciding according to which rule the approxi- mite polynomials, a particular case of which is the Tay-
mating function cp(t)=<P(f, t) should be chosen from the lor polynomial; in this case the values of a function
set W (when chosing the method of approximation) it is and of its first n consecutive derivatives at a point are
evident that one should strive for the highest possible interpolated by an algebraic polynomial of degree n.
accuracy of approximation and at the same time for Interpolation with splines has its peculiarities con-
simplicity in the construction of <P(f, t), taking account nected with the choice of the interpolation nodes and
of the information available about f The first require- with the boundary conditions ensuring the existence
ment leads to a function </>j in W that is 'nearest' to f, and uniqueness of the interpolating spline. For
i.e. </>j is such that instance, a spline s of degree r;:.2 and defect 1 with
pif, </>f) = inf pif, </. respect to the subdivision (1) taking prescribed values
,/>E!fI
at N different points t; of the interval (a, b) such that
Here immediately questions arise about the existence 'T;-I~ti~'T;, i=l, ... ,N, exists and is unique if the
and the uniqueness of such a function </>j (the function boundary conditions are given in terms of ma condi-
of best approximation), and also about its characteristic tions s(p)(a), O";;;;;p~r-l, and mb conditions s{Jt)(b),
properties (see [5]). Existence is ensured if W is a closed O";;;;;JJ.~r -1, with ma +mb =r. Functions
locally compact set, and in particular if W is a finite- fEG2r-k-l[a, b) can be put in one-to-one correspon-
dimensional subspace. Uniqueness depends both on the dence with splines s(j, t) of degree 2r - 1 and defect k,
properties of the approximating set W (see Haar condi- l.,;;;;;k ~r, with respect to the subdivision (1) by requir-
tion; Chebyshev system of functions) and on the metric ing that S(P)(j,'Ti)=/P)('Ti), p=O, ... ,k-l;
which defines piJ, </. There is a number of well-known i=l, ... ,N, where, if k<r, one also has to impose
necessary and sufficient conditions which should be some boundary conditions. When k = r, this spline is
satisfied by the function </>j of best approximation in a called a Hermite or local spline, because the behaviour
variety of cases (see Polynomial of best approximation; of the spline in an interval ('Ti -I, 'T;) is determined by
Chebyshev theorem; Markov criterion). However, as a the values of the function f and of its derivatives f(p) ,
rule, these criteria do not yield effective algorithms to p= 1, ... ,k -1, at the points 'T; -I and'Ti.
construct </>/t). Hence it is of great importance to have In the approximation of functions an important role
methods available which, on the basis of the informa- is also played by linear methods constructed on the
tion about the function f to be approximated, allow basis of an expansion of the function in a Fourier series
one to construct effectively some function <P(f, t) in W with respect to an orthogonal system. In particular, in
yielding an acceptable approximation. Here, in the the periodic case a widely used approximation tool con-
first place, one ought to mention linear methods (these sists of Fourier sums, and their various means (the case
satisfy cp(aJ/1 + a2/2 , t)=al</>(fh t)+a2</>(f2, t, to of the trigonometric orthogonal system) (see Approxi-
which, in particular, interpolation methods belong: hav- mation of functions, linear methods).
ing prescribed the points tl, ... ,tN in Q, one chooses The investigation and estimation of the error of
<P(f, t) from those functions cp in ~ that satisfy the approximation is important from the practical point of
interpolation condition view and at the same time is the branch of the approxi-
(2) mation of functions that is richest in ideas. More pre-

220

cisely, the development of methods dealing with
estimating the error, the investigation of this depen-
dence on the smoothness of the approximated function,
as well as the study and comparison of the approxima-
tion properties of various approximation methods, have
led to the creation of the approximation theory of func-
tions, one of the most rapidly developing branches of
mathematical analysis.
The foundation of the approximation theory of func-
tions was laid in the papers of P.L. Chebyshev in 1854
- 1859 (see [1)) on best uniform approximation of con-
tinuous functions by polynomials and rational fractions
and in the work of K. Weierstrass [2], who proved in
1885 that with any continuous function on an interval
[a, b], or on the whole real axis and having period 277,
there can be associated a sequence of algebraic (respec-
tively, trigonometric) polynomials Pn(j, t) of degree
n = 1, 2, ... , such that, as n-'HX),
where Q is [a, b], respectively the whole real axis. Simi-
lar assertions hold in the case when the measure of the
error of approximation is defined by an integral metric,
as well as in the case of functions of several variables.
Particular importance has been given to the investiga-
tion of the speed with which the sequence pf.j, Pn(j))
decreases in dependence on the properties of the
approximated function and the choice of the approxi-
mating polynomials Pn(j, t). The most attention has
been paid to the study of best approximation problems,
and also to that of linear methods of approximation
with the property that for a given f the polynomial
Pn(j, t) can be constructed effectively. An important
stage in the development of the approximation theory
of functions is connected with the names of Ch.J. de la
Vallee-Poussin, D. Jackson, and S.N. Bernstein [S.N.
Bernshteln], who initiated the investigation of the rela-
tion between the speed of decrease of the approxima-
tion when approximating the function f by means of
suitably chosen polynomials Pn(j, t) of degree n (as
n ~ (0), and the difference-differential properties of f It
turns out that in many cases these properties of f, e.g.,
f having derivatives, smoothness of those derivatives,
etc., can be characterized by the sequence of approxi-
mating polynomials and by the behaviour of the
corresponding error of approximation (see Approxima-
tion of functions, direct and inverse theorems). This
yields a new, constructive characterization of continu-
ous and differentiable functions. In the first third of the
twentieth century this kind of problem was dominating
approximation theory; for this reason this field was
also thought of as the constructive theory of functions.
In the nineteen-thirties and -fourties there appeared
papers by A.N. Kolmogorov, J. Favard and S.M.
APPROXIMATION OF FUNCTIONS
Nikol'skll which initiated a new direction of research
connected with the approximation of classes of func-
tions by finite-dimensional subspaces and with obtain-
ing accurate estimates of the error, using difference-
differential properties defining the class. The objective
is to obtain the quantities
sup p.(j, PN(f)),
where p.(g, cp) is the chosen measure of the approxima-
tion error, un is a class of functions and Pn(j, t) is the
approximating (generalized) polynomial whose coeffi-
cients are determined by the method of approximation.
Results of this kind make it possible to compare
approximation methods from the point of view of their
approximating capabilities and to formulate the prob-
lem, important in applications, of determining for a
given class of functions the optimal (best) approximat-
ing set (of fixed dimension N). Research in this direc-
tion, based both on the properties of specific approxi-
mation methods and on advanced results of functional
analysis, turned out to be also very fertile in ideas,
because, by proceeding in this way, new facts about the
mutual connections between many extremal problems
of various types were discovered, making it possible to
obtain deep and delicate relations in the theory of func-
tions. Owing to that it was possible to conclusively
solve a number of extremal problems concerning the
best approximation of important classes of functions
(see [5] and [7], and also Approximation of functions,
extremal problems in function classes).
Some other aspects of the approximation of functions.
Additional restrictions can be put on the approximat-
ing function cp(t)=cp(j, t) in 91. If they are not con-
nected with the function f (for instance, restrictions on,
or relations between, the coefficients of the approximat-
ing polynomial), then it is really a matter of determin-
ing 91 more precise. A new situation arises if the restric-
tions on cp(j, t) are connected with the approximated
function f; one of the interesting cases is that of one-
sided approximation. Here cp(j, t) in 91 is required to
satisfy the inequality cp(j, I)"';;'f(t) (or ~) and the error
is estimated in an integral metric (see [19)).
In applied problems, besides the case of explicity
given functions, one may have to approximate curves
and surfaces which are given only in parametric form;
as tools of approximation one can use, for instance,
parametric splines. It is most natural to measure the
error of approximation by means of the Hausdorff dis-
tance, which properly takes into account the geometric
proximity of such objects. For instance, for two curves
I} and 12 , the Hausdorff distance is defined by
r(/l, [2) = max {max min P(P, Q), maxminp(P, Q)},
Pe/l Qe/, Pel, Qe/l
221
APPROXIMATION OF FUNCTIONS
where p(P, Q) is the Euclidean (or any other) distance
between the points P and Q. Here the Hausdorff dis-
tance is preferable when choosing the measure of the
error of approximation; this is also the case in some
cases of the approximation of functions, for instance
when a discontinuous function has to be approximated
by a smooth function (see [16]).
The solution of a number of problems in the theory
of approximation of functions is closely connected with
the study of extremal properties of polynomials in one
or another concrete system of functions (inequalities for
the derivatives of polynomials, polynomials least deviat-
ing from zero, etc.). In particular, the proofs of inverse
theorems in the theory of approximation of functions
strongly rely on inequalities giving estimates of the
norm (or value at a fixed point) of some derivatives of
an algebraic or trigonometric polynomial, using certain
characteristics of the polynomial itself. In this direction
a number of precise results is known, and these have
generalizations to entire functions (see [6] - [10]).
The problem concerning an algebraic polynomial
(with prescribed leading coefficient) least deviating
from zero in the metric of C or 4 on the interval [a, b],
and the equivalent problem concerning the best
approximation of the function t n by a polynomial of
degree at most n -1, was investigated by Chebyshev
(for C) and by his students (for Ld. The solution is
given by the Chebyshev polynomials of the first kind
(C) and the second kind (L 1) and by the Legendre
polynomials (L 2 ), which have wide applications both in
theoretical and applied investigations. A number of
results is known for the more general case when all
coefficients of the polynomial are subjected to some
restrictions (see [6]). The problem concerning a mono-
spline of minimal norm, equivalent to that of finding
the best approximation of the function t n by splines of
degree n - 1 with free knots, is of special importance in
view of the fact that in a number of cases the problem
of obtaining the best quadrature formula reduces to it.
For the approximation of functions in the complex
plane see Approximation of functions of a complex vari-
able.
References
[I] CHEBYSHEV, P.L.: 'Problems about smallest quantities con-
nected with an approximate representation of a function
(1859)', in Complete collected works, Moscow-Leningrad, 1947,
pp. 151-235 (in Russian).
[2A] WEIERSTRASS, K.: 'Ueber die analytische Darstellbarkeit
sogenannter willkiirlicher Funktionen einer reellen
Veranderlichen', Sitzungsberichte Akad. Berlin (1885), 633-639.
[2B] WEIERSTRASS, K.: 'Ueber die analytische Darstellbarkeit
sogenannter willkiirlicher Funktionen einer reellen
Veranderlichen', Sitzungsberichte Akad. Berlin (1885), 789-805.
[3] GoNCHAROV, V.L.: The theory of interpolation and approxima
tion offunctions, Moscow, 1954 (in Russian).
[4] NATANSON, I.P.: Constructive function theory, 1-3, F. Ungar,
1964-1965 (translated from the Russian).
222
[5] KORNEiCHUK, I.P.: ExtTem(ll problems in approximation theory,
Moscow, 1976 (in Russian).
[6] DZYADYK, V.K.: Introduction to the theory of uniform approxi-
mation offunctiOns by polynomials, Moscow, 1977 (in Russian).
[7] TnrnOMIROV, V.M.: Some questions in approximation theory,
Moscow, 1976 (in Russian).
[8] NIKOL'sKIi, S.M.: Approximation offunctions of several variables
and imbedding theorems, Springer, 1975 (translated from the
Russian).
[9] ACHIEZER, N.I. [N.1. AKHIEzER]: Theory of approximation, F.
Ungar, 1956 (translated from the Russian).
[10] TIMAN, AF.: Theory of approximation offunctions of a real
variable, Pergamon, 1963 (translated from the Russian).
[11] KOROVKIN, P.P.: Linear operators and approximation theory,
Hindushtan Publ. Comp., Delhi, 1960 (translated from the
Russian).
[12] AHLBERG, I.H., NILSON, liN. and WALSH, I.F.: The theory of
splines and their applications, Acad. Press, 1967.
[13] STECHIUN, S.B. and SUBBOTIN, Yu.N.: Splines in numerical
mathematics, Moscow, 1976 (in Russian).
[14] LAURENT, P.I.: Approximation et optimisation, Herman, 1972.
[15] CoLLATZ, L. and KRAss, W.: Approximationstheorie: Tscheby-
scheffiche Approximationen mit Anwendungen, Teubner, 1973.
[16] SENDOV, B.: Hausdorff approximations, Sofia, 1979.
[17] NIKOL'sKIi, S.M.: Quadratureformulas, Moscow, 1979 (in Rus-
sian).
[18] ZAv'YADov, Yu.S., KVASOV, B.1. and MlROSHNICHENKO, V.I.:
Methods of spline functiOns, Moscow, 1980.
[19] KoRNEicHUK, N.P., LIGUN, AA and DoRONIN, V.G.: Approx-
imation with constraints, Kiev, 1982 (in Russian).
[20] DAVIS, P.J.: Interpolation and approximation, Dover, reprint,
1963.
[21] LoRENTZ, G.G.: Approximation offunctions, Holt, Rinehart
and Winston, 1966.
[22] KORNEicHUK, N.P.: Splines in approximation theory, Moscow,
1984 (in Russian).
N.P. Korneichuk
Editorial comments.
References
[A1] TODD, 1.: 'Applications of transformation theory; a legacy
from Zolotarev (1847-1878)" in S.P. Singh, J.H.w. Burry and
B. Watson (eels.): Approximation theory and spline functions,
Reidel, 1984, pp. 207-245.
[A2] SZEG<), G.: 'The contributions of L. Fejer to the constructive
function theory', in G. Alexits and S.B. Stechkin (eels.): Con-
structive theory of functions, Akad. Kiadn6, Budapest, 1972,
pp. 19-26.
[A3] IBRAGIMOV, 1.1.: 'On S.N. Bernstein's publications in the con-
structive theory of functions', in G. Alexits and S.B. Stech kin
(eels.): Constructive theory of functions, Akad. Kiadn6,
Budapest, 1972, pp. 27-40 (in Russian).
[A4] MICHELLI, C.A, PAl, D.V. and LIMAYE, B.V. (EDS.): Methods
of functional analYSis in approximation theory, Birkhauser,
1986.
[A5] CHENEY, E.W.: Introduction to approximation theory, Chel-
sea, reprint, 1982.
[A6] SCHUMAKER, L.L.: Spline functions, basic theory, Wiley,
1981.
[A7] RlvLIN, T.J.: The Chebyshev polynomials, Wiley, 1974.
[A8] DEVORE, R.A.: 'Freud's work in constructive function theory'
J. Approximation theory 46 (1986), 32-37.
[A9] TuRAN, P.: 'On some open problems of approximation
theory'. J. Approximation theory 29 (1980), 23-85. (Mat.
Lapok25 (1974),21-75)
[A 10] BRAESS, D.: Nonlinear approximation theory, Springer, 1986.
AMS 1980 Subject Classification: 41-XX, 41 A65,
41A10, 41A15

APPROXIMATION OF FUNCTIONS, DIRECT AND
INVERSE THEOREMS - Theorems and equalities
establishing a connection between the difference-
differential properties of the function to be approxi-
mated and the magnitude (and behaviour) of the error
of approximation, by various methods. Direct theorems
give an estimate of the error of approximation of a
function in terms of its smoothness properties (the
existence of derivatives of a given order, the modulus of
continuity of f or of some of its derivatives, etc.). In
the case of best approximation by polynomials, direct
theorems are also known as Jackson-type theorems [1],
together with their many generalizations and refine-
ments (see Jackson inequality and Jackson theorem).
Inverse theorems characterize difference-differential
properties of functions depending on the rapidity with
which the errors of best, or any other, approximations
tend to zero. The problem of obtaining inverse
theorems in the approximation of functions was first
stated, and in some cases solved, by S.N. Bernstein
[S.N. Bernshtei'n] [2]. A comparison of direct and
inverse theorems allows one sometimes to characterize
completely the class of functions having specific
smoothness properties using, for instance, sequences of
best approximations.
The connection between direct and inverse theorems
is most simple in the periodic case. Let C be the space
of 2'17"-periodic continuous functions on the whole real
axis with norm
I J I c= max
t
IJ(t) I;
let
be the best approximation of a function f in C by the
subspace Tn of trigonometric polynomials of degree at
most n, !et w(f, c5) be the modulus of continuity of f:
and let Cr , r = 1, 2, ... , be the set of functions in C
(;0 = C) that are r times continuously differentiable .sm
the whole real axis. A direct theorem states: If f E er,
then
n = 1,2, . , ., r = 0, I, ' .. ,
where the constant M does not depend on n. A
stronger assertion is: It is possible to indicate a
sequence of linear methods Un, n =0, 1, ... , associat-
ing a polynomial Un(f, DE Tn to a function f(t) in C
and such that for fEcr the error Ilf-Un(f) lie is
bounded by the right-hapd side of (1). An inverse
theorem states that for f E C
1IJ..61
w(f, 8) ,;;;; M8 Z E(f, Tn -d, 8>0,
n =1
where M is an absolute constant and [1/ c5]
APPROXIMA nOl\i Of IT;-";CTIOI\S, DIRECT Al\iD Il\iVERSE THEOREMS
integer part of 1 / c5. And if for some positive integer r
the series
n =1
converges, then it follows that fEcr and that, analo-
gous to (2), w(f(r), 1 / n) can be estimated in terms of
E(f, Tn -I), n = 1, 2, . .. (see [4], [8] and [12]). These
estimates imply, in particular, that if
E(f, Tn-I) = O(n- r - a ), r=O, I, ... , O<a<l,
then f E Cr and the fr)(t) satisfy for O<a< 1 the
Holder inequality
IJ(r)(t')-J(r)(t'') I,;;;; K I t'-t" 1\ (3)
whereas for a = 1 they satisfy the Zygmund inequality
/j<r V )-2J (r) [t' ~t"] +J(r)(t'') I ,; ; Kit' -t" I. (4)
Denoting this class of functions by Kir+ a , one obtains
a constructive characteristic for it: fEKHr+a if and
only if
A 2'17"-periodic function is infinitely differentiable on the
whole real axis if and only if
lim n rE(f, Tn -d =
n--+co
for all r= 1,2, ....
Similar properties hold for the approximation of
periodic functions in the metric of Lp(0,2'17")
(l";;;;p < 00), and also for (not necessarily periodic) func-
tions defined on the whole real axis when being
approximated by entire functions of finite order (see [7]
and_ [8]). Direct and inverse theorems are known for
fEer; as a difference-differential characteristic the
modulus of smoothness Wk(f, c5) of order k = 1,2, ... ,
of f (or of some of its derivatives) is used (see [4] and
[8]).
The situation is different in the case of approxima-
tion on a finite interval. Let C = C[a, b) be the space
of continuous functions on [a, b), with norm
II J lie = max IJ(t) I
a.;;;t.;;;b
(1)
Let er = er[a, b) be the set of the functions that are r
times continuously differentiable on [a, b), CO = C, and
let KHr+a[a, b) be the class of functions defined by the
inequalities (3) and (4) with t', t" E[a, b). For the best
approximation,
E(f,An-l;a,b) = pEA,,_1
inf IIJ-p lie, n=I,2, ... ,
of a function f E er by the subspace An - 1 of algebraic
polynomials of degree at most n - lone has an esti-
mate of the form (1) in terms of the modulus of con-
tinuity of f r ) on [a, b), but the converse, analogous to
(2)
that of the periodic case (with an inequality of the form
IS the (2)), now only applies for intervals contained in (a, b).
223
APPROXIMATION OF FUNCTIONS, DIRECT AND INVERSE THEOREMS
For instance, if
E(f,A n - 1 ;a,b) 0;;;; Mn- r -
r = 0, 1, ... , 0 < a < 1,
then one can only assert that f belongs to
KHr+"[a I, b d defined by the inequalities (3) (with
O<u<l) merely on the interval [aJ,bdC(a,b); the
constant K depends on a, ai, b 1 ,b, and can increase
unboundedly as a 1 ~a and b 1 ~b. There exist functions
not belonging to KHr+"[a, b] for which nevertheless
E(f,A n - 1 ;a,b) = O(n- r -
For instance
E(VI=t2,A n - 1; -1, 1) <~, n=I,2, ... ,
'/Tn
although VI=f2 fKH" on [-1, 1] for every u> 1/2.
It turned out that algebraic polynomials, retaining on
the whole interval [a, b] the best order of approxima-
tion of a function fEC, can yield at the end points of
the interval a substantially better approximation. This
phenomenon was first discovered by S.M. Nikol'skii
(see [3]). If, in particular, fEKHr+"[a, b], then for any
n >r there is a polynomial Pn(t)EAn -I such that
M[~v'(t-aXb-t)+ n\ r+
If(t)-Pn(t) I 0;;;;
a 0;;;; t 0;;;; b,
where the constant M depends neither on n nor on t.
Contrary to (5), this assertion has a converse: If for an
fEC there exists a sequence of polynomials Pn(t)EA n - 1
such that for some r =0, 1, ... , and O<u< 1 (6) is
satisfied, then fEKHr+"[a, b]. Direct and inverse
theorems for fEC[a, b] are known, stated in terms of
the modulus of continuity and moduli of smoothness
(see [4] and [8]).
Direct theorems giving order estimates for the error
of approximation in terms of difference-differential pro-
perties of the approximated function were established
for many concrete approximation methods (see [6], [8]
and [9], in particular for splines, splines of best approx-
imation and interpolation [10], cf. Spline).
Direct and inverse theorems for approximation in a
Hausdorff metric are known (see [13]). Some peculiari-
ties arise here; in particular, the characterization of
classes of functions by their best Hausdorff approxima-
tion is connected not only with the order of this
approximation but also with the magnitude of the con-
stant in the relevant inequality. For direct and inverse
theorems in the multi-dimensional case, see Approxima-
tion of functions of several real variables.
References
[I] JACKSON, D.: Ueber die Genauigkeit der Anniiherung stetiger
Funktionen durch ganze rationale Funktionen gegebenen Grades
und trigonometrische Summen gegebener Ordnung, Giittingen,
1911. Thesis.
224
[2] BERNSHTEiN, S.N.: 'On the best approximation of continuous
functions by polynomials of a given degree (1912)" in Col-
a, (5)
lected works, Vol. 1, Moscow, 1952, pp. 11-104 (in Russian).
[3] NIKOL'SKII, S.M.: Izv. Akad. Nauk SSSR Ser. Mat. 10, no. 4
(1946),295-317.
[4] DzYADYK, V.K.: Introduction to the theory of uniform approxi-
mation offunctions by polynOmials, Moscow, 1977 (in Russian).
[5] KORNEicHUK, N.P.: Extremal problems in approximation theory,
Moscow, 1976 (in Russian).
[6] 'flxHOMIROV, V.M.: Some questions in approximation theory,
Moscow, 1976 (in Russian).
[7] ACHIEZER, N.J. [N.J. AKHIEZER]: Theory of approximation, F.
Ungar, 1956 (translated from the Russian).
a ). [8] TIMAN, A.F.: Theory of approximation offunctions of a real
variable, Pergamon, 1963 (translated from the Russian).
[9] KOROVKIN, P.P.: Linear operators and approximation theory,
Hindushtan Publ. Comp., Delhi, 1960 (translated from the
Russian).
[10] STECHKIN, S.B. and SUBBOTIN, Yu.N.: Splines in numerical
mathematics, Moscow, 1976 (in Russian).
[II] DAUGAVET, I.K.: Introduction to the theory of approximation of
functiOns, Leningrad, 1977 (in Russian).
[I2] STECHKIN, S.A.: Izv. Akad. Nauk SSSR Ser., no. 3 (1951),
217-242.
[13] SENDOV, B.: Hausdorff approximations, Sofia, 1979.
[14] LORENTZ, G.G.: Approximation offunctions, Holt, Rinehart
and Winston, 1966.
N.P. Kornelchuk
Editorial comments.
References
a
, (6) [A 1] NATANsON, I.P.: Constructive theory of functions, 1-3, F.
Ungar, 1964-1965 (translated from the Russian).
[A2] DEVORE, R.A.: The approximation of continuous functions by
positive linear operators, Springer, 1972.
AMS 1980 Subject Classification: 41 A27, 41 A25
APPROXIMATION OF FUNCTIONS, EXTREMAL
PROBLEMS IN FUNCTION CLASSES - Problems con-
nected with the search for upper bounds of approxima-
tion errors with respect to a fixed class of functions and
with the choice of an approximation tool that is best in
some sense or other. The foundations of the research
concerning this kind of problems were laid by A.N.
Kolmogorov ([1], [2]), J. Favard ([3], [4]) and S.M.
Nikol'skii ([5], [6]). These researches were greatly
extended in the nineteen-fifties. They stimulated the
requirements of numerical mathematics, penetrating
ever deeper into problems of an optimization character.
If in a normed function space X one considers the
approximation of functions from a class we by func-
tions from a fixed set 91cX, then problems on deter-
mining the quantity
E(I]](,9C)x = supE(j, 9C)x, (1)
jE'lJI
where
E(j, \}c)x inf II f -<I> I x
= .pE'l/
is the best approximation of a function f(t) by 91, as
well as on determining the quantity
6' (I]]e, \}c, U)x sup II f -
= jdJl Uf II x, (2)
where U is a specific approximation method, given by
 APPROXIMATION OF FUNCTIONS, EXTREMAL PROBLEMS IN FUNCTION CLASSES
some operator acting from X into 91, are of interest.
Geometrically, the supremum in (1) characterizes the
deviation of the set we from 91 in the metric of X. In
practice, E(we,91)x can be looked upon, first, as the
minimal-possible bound from above for the best
approximation of 91 by functions that are only known
to belong to we, and secondly, it gives a tool for
measuring and comparing the approximative properties
of concrete approximation methods on we. In (2), the
most important case is when 91 = 91N is an N-
dimensional space and U is a linear approximation
method. A whole series of exact and asymptotically-
exact results is known for the approximation of func-
tion classes by concrete linear methods (in particular,
by polynomial or spline methods), see [l] - [12], [19].
However, special interest attaches to methods attaining
the infimum
com, inNh = UXC'1I
inf C([Q, inN, U)x,
N
over all linear bounded operators U from X into 91x ,
i.e. to all linear methods that are optimal for we. Obvi-
ously,
and there naturally arises the question on whether the
equality sign can be attained. Apart from the trivial
case when X is a Hilbert function space and the best
approximation of a function is given by its Fourier
sums with respect to an orthogonal basis of 91 x , some
results when linear methods realize the best
approximation on the entire class we are also known for
non-Hilbert spaces.
Thus, if X is the space L = Lp[O, 2'17'] (I';;;;;p,;;;;; 00) of
2'IT-periodic functions, 91~ _I is the subspace of tri-
gonometric polynomials of order
(dim 91fn-1 =2n-l),_and if MW;, r=I,2, ... , is the
class of functions IELp for which I(r-I)(t) is absolutely
continuous on [0,2'17'] and for which 11/(r) IIi, ,;;;;;M, then
E(MW;, inf,,-I)i, = C(MW;, inf,,-di, = MKrn- r,
p = 1,00; n, r = 1,2, ... ,
where Kr is_ the Favar~ constant. The best approxima-
tion on MWl and MW~ is obtained, moreover, by a
linear method U~ -I, constructed on the basis of
Fourier sums (cf. Approximation of functions, linear
methods, formula (3 for a specific choice of the multi-
plication coefficients A~n -I) =A~n -I)(r). Linear opera-
tors, with values in 91~ -I, yielding the supremum of
the best approximatiol! on convolu!ion classes, includ-
ing, in particular, MW~ and MWl with a rational
number r>O, as well as on classes of conjugate func-
tions, have been established (cf. [10], [11]).
For approximation by the subspace Sq;, of 2'17'-
periodic splines of order m and defect 1 with knots
(nodes, joints) at k'IT / n (dimS~ =2n), the equalities
E(Mwr sr-I),
P' 2n L" = C(MwrP' sr-I),
2n L"
= MKn- r
r'
p = 1,00; n, r = 1,2, ... ,
are valid. The splines (Jr-I (j, t) in S~; 1 interpolating a
function I at k'IT / n if r is even, and at (2k + I)'IT / 2n if
r is odd, provide the best linear methods in this case.
Relative to MW~ these splines have exceptional
approximative properties, _ since they_ give the best
approximation of an IEMW~ in any Lp (cf. [7], [15]).
When (1) and (3) coincide and when it is possible to
construct a concrete linear method solving both prob-
lems, the cases considered are, in a well-known sense,
optimal. In other cases it proved expedient in solving
(1) to adopt an approach based on general duality
theorems, reflecting fundamental relations in geometric
convex analysis (see [7], [8]). If, e.g. X is an arbitrary
linear normed space, X is its dual and 91 is a convex
(3)
set in X, then for any x EX
inf
ue91
II x-u II = sun
Fet,
[F(X)-SUPF(u)];
UE~
(4)
II F II..;I
in particular, if 91 is a subspace,
inf II x-u I = SUp {F(x): FEX, II F 11';;;;1, (5)
ue'1l
F(u)=O for all uEin}.
In a number of cases (4), or (5), allows one to reduce
the calculation of the supremum (1) to the more tran-
sparent problem of determining the extremum of an
explicitly defined functional on some set of functions,
which, if 91 is a subspace, is related to orthogonality
conditions. For example, by (5) the estimation of
E(Wq, 91fn-I)i" I~p, q';;;;;oo, can be reduced, using
known inequalities, to the calculation of the supremum
n- 1
of the norms I g Ilq', q' =q / (q -1), on the set of func-
tions gE W;', p' =p / (p -1), for which
2.. coskt
[g(t)Sinktdt = 0, k=O, ... ,n-l.
More refined situations occur when (1) has to be solved
on classes not given by restrictions on the norm of the
r-th derivative j<r)(t), but on its modulus of continuity
w(j(r) .z.~)x. Particularly int~resting the case where is
we=W'H'" (r=O,I, ... ; WOH"'=H"') is the class of
c
2'IT-periodic functions IE r (CO = C) for which
w(j(r),~)c = w(j(r),~) .;;;; w(~),
where w(~) is a given modulus of continuity, e.g.
w(~)=M~ (0 < a';;;;; 1, O';;;;;~';;;;;'IT). Here the application of
(5) requires the use of fine results on differentiable
periodic functions of the type of the comparison
theorem and the Kolmogorov inequality (for the norm
of derivatives), here described using the tool of rear-
rangements (of equi-measurable functions). If w(~) is
convex from above, the following equalities ([7], [15])
225
APPROXIMA T10N OF FUNCTIONS, EXTREMAL PROBLEMS IN FUNCTION CLASSES
are valid:
E(WrH"', ~L-I)x = E(WrH"', S5,,)X = lIf"r(CAl)ilx,
n = 1,2, ... ; r = 0, 1, ....
- -
Here X=C or LJ, thef"r(w,t) are functions in W"Hw
of period 2'1T / n, with zero average over a period, for
which fi,~) is even, equal to [w('1T /n-2t)]/2 on
[0, '1T / 2n], and equal to - [.>(2t -'1T / n)1/ 2 on
['1T / 2n, '1T / n]. The norms II f"r(w) II x have an explicit
expression, e.g.
II f" 0 (CAl) lie = ~ CAl [: J.
IIf"r(CAl) lie = 2!r [~r('/T-t)CAl [~) dt,
where the functions ~k(t) on [0, '1T] are given recurrently
by ,,_/
~I(t) = t, ~k(t) = ~ I ~k-I(U)du.
The best linea.,r method from Cinto 91~-1' or into S~n'
for the class WH W is known for w(8)=m8, 0~8~'1T, i.e.
when Wr H W =MW"oo+ 1 . The interpolation splines
or(j, t) in S~'IT yield the supremum E(W" H W , S~'IT)c (for
any convex w(8 only if r = 1.
For solutions to extremal problems for function
classes on a finite interval and not restricted by rigid
boundary conditions, it is not possible to give such
complete results as in the periodic case; the extremal
end points of the interval have a perturbing influence
on the extremal functions, which require an increase in
the order of differentiability. Here, some exact asymp-
totic results are known. If MWr Ha, r =0, 1, ... ,
O<a<l, MWOHa=MH a, is the class of functions
fEel -1, 1] for which
IprV)-j<r)(t'') I.;;; M I t'-( la
then for the best uniform approximation on [- 1, 1] by
the subspace 91~ of algebraic polynomials of degree
n - 1, the relations
1T
lim nr+aE(MwrH a, ~~)c
n_oo
= M2 f~r('IT-t)tadt,
0
r = 1,2, ... ,
hold. It is useful to compare these results with the
corresponding ones in the periodic case. For a = 1, the
right-hand sides of (6), (7) are equal, respectively, to
MK I and MKr + I. Dropping the requirement of a
polynomial of best approximation, one obtains stronger
results, essentially improving the approximation at the
end points of [-1, 1] without losing the best asymptot-
ics on the whole interval. E.g., for any fEMHa there is
a sequence of algebraic polynomials Pn(j, t)E91: such
226
that uniformly in tEl -1,1], as n~oo, one has
I j(t)-Pn(j, t) I .;;; ~ [[ : v'l=t2r +o(n- a)]
- = E(MHQ, ~:)d(l-t2)a/2 +0(1)].
Analogous results hold for MW" HI, r= 1, 2, ... , see
[11 ]. In spline-approximation problems (both of best
and of interpolation type) certain exact and
asymptotically-exact results are known for function
classes on intervals (cf. [15]).
In the case of one-sided approximation (in an
integral metric) a series of exact results concerning the
estimation of the error of best approximation by poly-
r=I,2, ... , nomials and splines regarding the above-mentioned
function classes is known (cf. [14]). In deriving them,
essential use is made of duality relations for the best
approximation under restrictions given by cones.
The search for a best approximation tool (of fixed
dimension), leads, for a fixed function class W1, to the
problem of widths: Determine (cf. (1), (3:
dN(IJR, X) = infE(IJR, ~N)x'
!liN
where the infima are taken over all subspaces 91N of X
(as well as over their translations) of dimension N, and
also to determine extremal (best) subspaces yielding
these infima. An upper bound for dN and d~ is given
by E(W1,91)x and C(W1,91N )x, which have been esta-
blished for concrete subspaces 91N . The fundamental
difficulty in the problem of widths, here, is to obtain
exact lower bounds. In a number of cases it is possible
to obtain such bounds by appealing to topological con-
siderations, in particular, Borsuk's antipodal theorem
(cf. [8]). In practically all cases of exactly solving best
(t',t"E[-I,I]),
approximation problems in the classes MW; and
W" H W of periodic functions by the subspaces 91fn-1
(trigonometric polynomials of order n - 1) or S~n
(splines of order m and defect 1 with respect to the
partition of knots k'1T / n), the known exact upper
(6)
bounds E(W1, 91N )x give also the values of dN for these
classes. Moreover, for periodic classes d 2n - 1 =d2n In
(7) particular (cf. [7], [8], [15], [16])
d 2n - UV';", C) =
1 d 2n (J.v';", C) = d 2n - UV}, Ld =
1
= d 2n (W},L 1) = Krn- r, n,r=I,2, ... ,
and for an upper convex w(8) and for X=C or LIane
has
d2n -
-
1(W'Hw,X)
-
= d2n (W'HW,X) = 11f",(CAl)llx,
n=I,2, ... ; r=O, 1, ....
It should be observed that \)(~ -I is extremal for the
classes considered, for all r, and that no subspace of
dimension 2n gives a better approximation for these
 APPROXIMATION OF FUNCTIONS, EXTREMAL PROBLEMS IN FUNCTION CLASSES
classes than 9Cf" -I (which is of dimension 2n - 1). The
s~bspace ~Tn of splines is ~xtremal for W~+ I and
W' H W in C if m =r ~d for _W;+I ~ LI if_ m~r. The
linear widths d~ of W~ in C and W~ in LI coincide
with dN , they are attained on 9Crn-I and Ss.;1 by the
best linear methods that were referred to above. The
widths dN and d~ of Ws. in L2 for N=2n -1 and
N = 2n are equal and are obtained by Fourier tri-
gonometric sums. Exact asymptotic results for widths
of function classes on an interval are known in a
number of cases: The widths dN and d~ of W'oo in
C[ - 1, 1] coincide and are obtained for interpolation
splines of order r - 1 with respect to non-uniform parti-
tions (cf. [8]).
The problem of estimating the approximation error
on the set xr of r-fold integrals of functions from X
(which is not locally compact) can be stated properly if
for JEX r and fixed y>O one estimates
E(j, IflNh
w(j(r) , 'Yh
(w(g, o)x is the modulus of continuity of g in X) or (in
the case of a concrete approximation method) if one
estimates
Ilj- UN/llx
w(j(r), 'Yh .
Determining the suprema of these quantities on x r is
equivalent to the search for the least constant in the
corresponding Jackson inequality; subsequently one
may deal with the minimization over all subspaces of
dimension N. In a number of cases, these problems
have been solved. E.g. in the inequality
E(j, SS,)c ~ Mrn-rw ~(r), : t, jEer,
the least constant is Mr =Kr /2, and it cannot be
improved if sS.n is replaced by any other subspace of
the same dimension (d. [15]). Exact constants in the
Jackson inequalities for trigonometric approximation in
the uniform and in an integral metric are known (cf.
[7]). In the non-periodic case there are asymptotic
results available.
Among approximation problems in which the
approximating set is not a linear manifold but a convex
set, the problem of approximating one function class we
by another weI with best, in some sense or other,
smoothness properties, is of interest. Such a problem
arose initially at an intermediate stage in obtaining
exact bounds for
E(H"', Iflfn-l)C (g:Q = H"', g:Ql = MW),,)
(d. [7]); later it became to be considered as an indepen-
dent problem, and in a number of cases it was possible,
using (4), to obtain exact results.
For
g:Q = M1JV;,
g:Ql = M2W~ (O<r<k)
relations have been obtained between inequalities for
norms of derivatives and best approximations of a dif-
ferentiation operator by linear bounded operators (cf.
[13]).
Many extremal problems in the approximation of
functions may be interpreted as problems of optimal
recovery (cf. [15]). Let the information on a function
JEXbe given by a vector T(f, A)={Al(f), ... ,An(f)},
where Ak are given functionals on X (e.g. values of J
and/or some derivatives in fixed points). Given that J
belongs to a class we, it is required to reconstruct J or
some linear functional L(f) of it (e.g. L(f)= AI},
f b
L(t)= J(t)dt, etc.) on the basis of T(f, A) with least
possible error. The minimization need not only take
into account methods S juxtaposing T(f, A) with a
function q:,(t)~J(t) (or a functional 1(f)~L(f, but
also involves the choice of functionals Ab k = 1, ... ,N.
Depending on the choice of the error measure and the
class of methods S, the problem of optimal recovery of
a function may sometimes be reduced to determining
the widths dN , d~, the Chebyshev centre, or other
characteristics of we. The problem of optimal recovery
of f b
J(t) dt on the basis of information of the form
{{(tk)} or {{(v)(tk)} leads to the problem of best qua-
drature formulas for we. The best tool of recovering is
achieved, in a number of cases, by splines; e.g. the
splines (Jr-l (f, t) in ss.; 1 that inte_rpolate at equally-
distant points tk reconstruct an JE W~ on the basis of
information Ak(f) = J(td in each_ point t=/=tk with
minimum error on the whole class W~.
In spaces of functions of two or more variables,
excluding the trivial case of approximation in Hilbert
spaces, exact solutions of extremal problems have not
yet been obtained (1983). In a few cases an asymptotic
relation for the error of uniform approximation in func-
tion classes by Fourier sums or some average Fourier
sums is known (cf. [12]).
References
[I] KOLMOGOROY, A.N.: Ann. of Math. 36, no. 2 (1935), 521526.
[2] KOLMOGOROFF, A.N. [A.N. KOLMOGOROY]: 'Ueber die beste
Annliherung yon Funktionen einer gegebenen Funktionen-
klasse', Ann. of Math. 37, no. 1 (1936), 107-110.
[3] FAYARD, J.: C.R Acad. Sci. 203 (1936), 1122-1124.
[4A] FAVARD, 1.: 'Sur les meilleurs procedes d'approximation de
certaines classes de fonctions par des polynomes
trigonometriques', Bull. Sci. Math. 61 (1937), 209-224.
[4B] FAYARD, 1.: 'Sur les meilleurs procedes d'approximation de
certaines classes de fonctions par des polynomes
trigonometriques', Bull. Sci. Math. 61 (1937), 243-256.
[5] NIKoL'sKIi, S.M.: 'Approximations of periodic functions by
trigonometric polynomials', Trudy Mat. [nst. Steklov. 15 (1945),
1-76 (in Russian). English abstract.
227
APPROXIMATION OF FUNCTIONS, EXTREMAL PROBLEMS IN FUNCTION CLASSES
[6] NIKOL'sKIi, S.M.: 'Mean approximation of functions by tri-
gonometric polynomials', Izv. Akad. Nauk SSSR Ser. Mat. 10
(1946), 207-256 (in Russian).
[7] KORNEicHUK, N.P.: Extremal problems in approximation theory,
Moscow, 1976 (in Russian).
[8] TnrnOMIROV, V.M.: Some questions in approximation theory,
Moscow, 1976 (in Russian).
[9] DZYADYK, V.K.: Introduction to the theory of uniform approxi.
mation by polynomials, Moscow, 1977 (in Russian).
[10] ACHIEZER, N.!. [N.I. AKHIEZER]: Theory of approximation, F.
Ungar, 1956 (translated from the Russian).
[ II] nMAN, AF.: Theory of approximation offunctions of a real
variable, Pergamon Press, 1963 (translated from the Russian).
[12] STEPANETS, AI.: Uniform approximation by trigonometric poly
[13]
nomials. Linear methods, Kiev, 1984 (in Russian).
ARESTOV, V.V.: 'Some extremal problems for differentiable
functions of one variable', Trudy Mat. Inst. Steklov. 138 (1975),
3-28 (in Russian).
(14] KORNEicHUK, N.P., LIGUN, A.A and DORONIN, V.G.: Approx
imation with contraints, Kiev, 1982 (in Russian).
(15] KORNEicHUK, N.P.: Splines in approximation theory, Moscow,
1984 (in Russian).
N.P. Korneichuk
Editorial comments. The recent book by A. Pinkus [A1]
contains a comprehensive and valuable bibliography.
References
[A1] PINKUS, A: nwidths in approximation theory, Springer,
1985.
[A2] MICHELLI, C.A and RIVLIN, TJ.: 'A survey of optimal
recovery', in CA Michelli and T.J. Rivlin (eds.): Optimal esti-
mation in approximation theory, Plenum Press, 1977, pp. 1-
54.
AMS 1980 Subject Classification: 41 A65, 41 A50,
41A46
APPROXIMATION OF FUNCTIONS, LINEAR
MEfHODS - Approximation methods defined by linear
operators. If in a linear normed function space X a
linear manifold (linear subspace) 'in is chosen as
approximating set, then any linear operator U which
transforms a function fEX into a function
U(j, t)=(Uj)(t)E'in so that
U(adl +ad2' t) = al U(jJ, t)+a2U(h, t)
(where lXI and lX2 are arbitrary numbers), defines a
linear approximation method for functions in X by
means of functions in 'in. A linear approximation
method is called projective if U(j, t) = f(t) for all f in 'in;
it is called positive if U(j, t);;.O for non-negative func-
tions f
The finite-dimensional case is the most interesting
one. Here 'in = 'inN is an N-dimensional subspace, and
then
N
U(f, t) = UN(f, t) = ~ Ck(j)ct>k(t),
k =1
where {q:,k(t)}f is a basis of 'inN and Ck are certain
linear functionals defined on X. The choice of the
linearly independent system {q:,k(t)}f and the set of
functionals {Ck} f is determined by the information
about the function that one intends to use while con-
structing the linear method. If Ck(f) = f(tk), where
228
{tdf is a fixed system of points in the domain of
definition of f, and if q:,k(ti)=O when i=/=k, while
q:,k(tk) = 1, then UN(j, tk)= f(tk), k= 1, ... ,N, and one
has an interpolation method (for instance, Lagrange's
interpolation polynomial, or an interpolation spline). If
X = H is a Hilbert space of functions and Ck(f) are the
Fourier coefficients of a function f with respect to an
orthonormal system {q:,k(t)}, then the sum at the right-
hand side of (1) yields a linear method of orthogonal
projection of X onto 'inN; in this case
Ilf- UN(j)IIH = ~IJt-k~lak~kt
Thus, this method yields the best approximation of f by
linear combinations of the functions q:,k.
In the theory of linear approximation methods, spe-
cial attention is given to problems of convergence. Let
X be a Banach space, let {q:,1 (t), !f>z (t), . . .} be a
linearly independent system of functions on X, let 'inN
be the subspace generated by the first N functions of
this system, N = 1, 2, ... , and let UN be bounded
linear operators from X into 'inN, N= 1,2, . ... The
convergence UN(f, t)-'?f(t) (in the sense that
I UN - f II x-'?O as N -'? 00) for any function f in X
holds if and only if: 1) the sequence of norms I UN I
of the operators UN is bounded (see Banach-Steinhaus
theorem); and 2) UN(j, t)-'?f(t) for all functions fin a
set A that is everywhere dense in X. These conditions
are satisfie_d, ~. particular, in the space of 2'IT-periodic
functions Lp = Lp[O, 2'17], with 1<p < 00, and the opera-
tors Sn defined by the trigonometric Fourier sums
n = 1,2, ... ,
of f As the set A on which 2) is satisfied one can take
the set _of ~ trigonometric polynomials. If X is the
space C = qo, 2'17] (of continuous functions _on the
whole real axis with period 2'17) or if X equals L I, then
II Sn I ~lnn as n-'?ooJsee Lebt:sgue constants). Conse-
quently, there are in C and in LI functions f to which
the sequence {Sn(j, t)} does not converge in the
appropriate metric. Since in the Banach space of func-
tions X with a norm that is invariant with respect to
translations the linear operator U;, projecting X onto
the subspace Tn of trigonometric polynomials of degree
at most n satisfies the i~equa1itx. I U;, I ;;.11 Sn I (see
(I)
[3]), the divergence in C and L I also holds for the
sequence {U;'}. In particular, this is true in C for the
sequence of Lagrange interpolation operators with any
triangular matrix of interpolation nodes. Similar facts
of a negative character also hold in the non-periodic
case for the linear projection operators of the spaces
qa, b] and Lda, b] onto the subspaces An,

n = 1,2, ... , of algebraic polynomials of degree at
most n.
The diyergence of the sequence (2) for some func-
tions in C made it necessary to introduce and examine
various averages of Fourier sums that do not have this
disadvantage. Such are, for instance, the Fejer swnma-
tion method and the Bernstein - Rogosinski swnmation
method, which are particular cases (with a specific
choice of the numerical coefficients Ar) of polynomials
of the form
ao n
(3)
U~lf,t) = 2+ ~>W)(akcoskt+bksinkt),
k=1
n = 1,2, ....
Since
U~lf, t) =-
'Tr 0
f
1 2.". [ 1
-2 + ~>W) cosk(t -u) f(u)du,
n
k =1
]
(4)
the averages (3) belong to the very wide class of linear
approximation methods which can be represented in
the form of a convolution of fwith some (singular) ker-
nel whose properties (in this case, the properties of the
triangular matrix of numbers {Ar)}) determine the
answer to the problem of convergence. If
lim A\,") = 1, k=I,2, ... ,
n~oo
then the sums (3) converge uniformly when f is a tri-
gonometric polynomial, and so 2) is satisfied. For the
boun9edness of the norm of U~ as an operator from C
into C, it is necessary that
Aln) A~n)
-n-+ ... +-1- = 0(1)
and Ar) = 0(1), uniformly in n and k = 1, ... ,n.
These conditions are sufficient as well for the bounded-
ness of the sequence {II U~ II} if the matrix {A~n)} is
subject to some additional restrictions (for instance,
convexity or concavity of the columns). Similar to (3),
using the matrix of coefficients {A~n)} one can also con-
struct averages by means of Fourier - Chebyshev sums
(for fEC[ -1, I]) and also by means of the Lagrange
interpolation polynomials with nodes 2k7T / (2n + I) (in
the periodic case) or with nodes cos[(2k -1)7T / (2n + 2)]
on [ - I, 1] (see, for instance, [6]).
The question of convergence of positive linear opera-
tors U;; acting from C[a, b] into An or from Cinto Tn
(in particular, of operators of the form (4) with a posi-
tive kernel) is solved in terms of three test functions
(see [1]): For the uniform convergence of the sequence
U;;(f,t) to f(t)EC[a,b] or to f(t)EC it is necessary
and sufficient that this occurs for the functions 1, t and
t 2 , or, respectively, for the functions 1, sint, cost.
The investigation of the error of approximation by a
linear approximation method leads, mostly, to studies
on the rate of convergence of UN(f, t) to f(t), to esti-
APPROXI~IATION OF FUNCTIONS. LINEAR METHODS
mates of the error in terms of the difference-differential
properties of the approximated function, and to
answering the question of how the linear approximation
method behaves when the smoothness properties (of the
approximated function) are improved.
In the examination of approximating properties of
the linear approximation method (1), a natural role is
played by the best approximation Elf, ~N) of fby ele-
ments of the subspace ~N' The Lebesgue inequality
IIf-Sn(f) lie';;; Elf, Tn)C(lnn+3) (5)
shows, in connection with the Jackson inequality
v(r) , ! ]
Elf, Tn)e .;;; Mnw
(where wJg, 8) is the modulus of continuity of a func-
tion g E C), that, although the order of approximation
by Fourier sums is slightly worse than the best approxi-
mation (In n in (5) cannot be replaced by a constant),
these sums are affected by any increase of the order of
differentiability of the function to be approximated.
For some linear approximation methods the order of
approximation cannot be higher than a certain quan-
tity, no matter how many derivatives the function f
may have (the phenomenon of saturation). Thus, the
order of approximation by positive linear polynomial
operators U;; cannot be higher then O(n -2); for Fejer
sums the saturation order is O(n -I); and for
Bernstein - Rogosinski sums it is O(n -2). Interpolation
splines with a specific choice of the knots (nodes,
joints) ensure the best order of approximation not only
for the function itself, but also for some of its deriva-
tives, depending on the degree of the polynomials out
of which the spline is composed (see [7], [8]).
In some particular cases, for specific linear approxi-
mation methods exact or asymptotically-exact estimates
for the errors of approximation were obtained for some
classes of functions. Strictly speaking, the first non-
trivial result of this kind was established by A.N. Kol-
mogorov, who proved in 1935 that
where W r K, r = 1, 2, ... , is the class of functions fEe
for which f(r-I)(t) is absolutely continuous on [0,27T],
and for which almost everywhere I f(r)(t) I,.;;;,K.
Subsequently, results of similar character were obtained
for Fourier sums (and some of their averages) with
respect to other important classes of functions, for
instance, in terms of upper bounds on the modulus of
continuity of the r-th derivative (see [2]), [6] and [9]).
Special interest is given to linear approximation
methods (I) yielding, on some class of functions, the
least upper bound of the best approximation by the
subspace ~N' Such a property for the classes W r K,
229
APPROXIMATION OF FUNCTIONS, LINEAR METHODS

r = I, 2, ... , with a specific choice of the A~n), applies the metric in a function space containing both f and <1>.
to sums of the form (3). For instance, for r = lone For example, if f and cp are continuous functions on a
should take segment [a, b], the uniform metric of C[a, b] is com-
A(n) hr k'TT
k = 2n+2 ctg 2n+2; monly used, i.e. one puts
the property also applies to interpolation splines of p.(j, q,) max I /(t)-<f>(t) I
= a.;;/.;;b
order r - I and defect I with knots hr / n, k = 0, I, ...
If continuity of the approximated function is not
(see [4] and also Approximation of functions, extremal
guaranteed or if the conditions of the problem imply
problems in function classes; Best linear method).
that it is important that f and <I> are close on [a, b] in
References
[IJ KOROYKIN, P.P.: Linear operators and approximation theory, an average sense, the integral metric of a space Ip[a, b]
Hindushtan Pub!. Comp., Delhi, 1960 (translated from the may be used, putting
Russian). b
[2J DZYADYK, V.K.: Introduction to the theory of uniform approxi.
mation offUnctions by polynomials, Moscow, 1977 (in Russian).
p.(j, q,) = f q(t) I/(t)-<f>(t) IP dt,
a
P >0,
[3J TnrnOMIROY, V.M.: Some questions in approximation theory,
Moscow, 1976 (in Russian). where q(t) is a weight function. The case p =2 is most
[4J KoltNEicHUK, N.P.: Extremal problems in approximation theory, often used and is most convenient from a practical
Moscow (in Russian). point of view (cf. Mean-square approximation of a
[5J GoNCHAROY, V.L.: The theory of interpolation and approxima
tion offUnctions, Moscow, 1954 (in Russian).
function).
[6J TIMAN, AF.: Theory of approximation offUnctions of a real The measure of approximation may take into
variable, Pergamon, 1963 (translated from the Russian). account only values of f and cp in discrete points tk,
[7J AHLBERG, l.R., NILSON, E. and WALSH, l.L.: The theory of
k=l, ... ,n, of [a, b], e.g.
splines and their applications, Acad. Press, 1967.
[8J STECHKIN, S.B. and SUBBOTIN, Yu.N.: Splines in numerical p.(j, q,) max I /(tk)-<f>(tk) I,
= l.;;k';;n
mathematics, Moscow, 1976 (in Russian).
[9J STEPANETS, AI.: Uniform approximation by trigonometric poly n
nomials. Linear methods, Kiev, 1981 (in Russian). p.(j, q,) = ~ qk I /(td-<f>(tk) IP,
[IOJ KORNEicHUK, N.P.: Splines in approximation theory, Moscow, k=l

1976 (in Russian). where qk are certain positive coefficients.

[11J RICE, l.R.: The approximation offunctions, I. Linear theory,
Addison-Wesley, 1964.
One defines in an analogous way the measure of
[12J ScHUMAKER, L.L.: Spline functions, basic theory, Wiley, 1981. approximation of functions in two or more variables.
[13J DAVIS, P.l.: Interpolation and approximation, Dover, reprint, The measure of approximation of a function f by a
1963.
N.P. Korneichuk family F of functions is usually taken to be the best
Editorial comments. The fact that three test functions approximation:
are enough for uniform convergence of a sequenc~ of posi- E(f, F) = p.(j, F) = .pEF
inf p.(j, q,) .
tive linear operators from C[a, b] into An (or from C into Tn)
can be found in, e.g., [A 1], Chapt. 3, Sect. 3. The theorem The quantity
is a generalization of the Bernstein theorem, and was
obtained by H. Bohman [A2]. For integral operators it was
E(Wl, F) = p.(Wl, F) = sup inf /L(f, cp)
fE~IcF

obtained by P.P. Korovkin (see also [1]).
References
[A1J CHENEY, E.W.: Introduction to approximation theory, Chel-
sea, reprint, 1982.
[A2J BOHMAN, R.: 'On approximation of continuous and of ana-
lytic functions', Arkiv Mat. 2 (1952), 43-56.
[A3] KrnSEWETIER, H.: Vorlesungen Dber lineare Approximation,
Deutsch. Verlag Wissenschaft., 1973.
[A4] DEVORE, R.A: The approximation of continuous functions by
positive linear operators, Springer, 1972.
[A5] CHENEY, E.W., HOBBY, G.R., MORRlS, P.D., ScHURER, F. and
WUBERT, D.E.: 'On the minimal property of the Fourier pro-
jection', Trans. Amer. Math. Soc. 143 (1969), 249-258.
is usually taken as the measure of approximation of a
class we of functions f by functions cp from a certain
fixed set F. It characterizes the maximal deviation of
functions in we from functions in F that are closest to
them.
In general, when approximation in an arbitrary
metric space X is considered, the measure of approxi-
mation JL{x, u) of an element x by an element u (a set
F) is the distance p(x, u) (or p(x, F) between x and u
(or F) in the metric of X.

AMS 1980 Subject Classification: 41 A36, 41 A45,

References
[I] GONCHAROY, V.L.: The theory of interpolation and approxima
41A40 tion offunctions, Moscow, 1954 (in Russian).
[2] NIKOL'SKIi, S.M.: Approximation offunctions of several variables
APPROXIMATION OF FUNCOONS, MEASURE and imbedding theorems, Springer, 1975 (translated from the
OF - A quantitative expression for the error of an Russian).
approximation. When the discussion is about the [3] RICE, J.R.: The approximation of functions, 1-2. Addison-
Wesley, 1964-1968.
approximation of a function f by a function <1>, the N.P. Kornetchuk
measure of approximation J1Ij, <1 is usually defined by v.P. Motornyt

230

Editorial comments. The measure of approximation is
also called the error measure.
References
[A1] LORENTZ, G.G.: Approximation of functions, Holt, Rinehart
and Winston, 1966.
[A2] PINKUS, A.: n-widths in approximation theory, Springer, 1985
(translated from the Russian).
AMS 1980 Subject Classification: 41 A50
APPROXIMATION OF FUNCTIONS OF A COM-
PLEX VARIABLE - The branch of complex analysis
studying problems regarding the approximate represen-
tation (approximation) of functions of a complex vari-
able by means of analytic functions of a specific class.
The fundamental problems in the theory of approxima-
tion of functions of a complex variable are: the possi-
bility of approximating; the rate of approximation; and
the approximation properties of various methods of
representing functions (interpolation sequences and
series, series in orthogonal polynomials or Faber poly-
nomials, expansion in continued fractions and Pade
approximation, sequences of polynomials in exponen-
tial functions, Dirichlet series, etc.). The theory of
approximation of functions of a complex variable is
intimately connected with other branches of complex
analysis, and with mathematics in general. Methods
and results on conformal mapping, integral representa-
tion, potential theory, the theory of function algebras,
etc., play an important role in approximation theory.
The central problems in the theory of approximation
of functions of a complex variable relate to approxima-
tion by polynomials or rational functions, in particular
by polynomials and rational functions of best approxi-
mation (existence, characteristic properties, uniqueness),
as well as to extremal problems and various estimates
for polynomials and rational functions (growth esti-
mates, inequalities for derivatives, polynomials and
rational functions, least deviation from zero, etc.).
A.A. Gonchar
The approximation of functions of a complex variable by
polynomials and rational functions. In this branch of
approximation theory one may distinguish several
directions.
I) The study of the possibility of approximating a
function f(z) of a complex variable z with given accu-
racy by polynomials or rational functions in z, in
dependence on the properties of the set E on which f is
given and on which f is to be approximated, on the
properties of the metric p of deviation and on the pro-
perties of f itself.
2) The study of the properties of polynomials and
rational functions of best approximation, i.e. polynomials
Pn(z;j, E, p) and rational functions Rn(z;j, E, p) of
degree not exceeding n, n =0, I, ... , for which
APPROXIMATION OF FUNCTIONS OF A COMPLEX VARIABLE
del
p(j, Pn(z;j, E, p == En(j, E, p) == inf{p(j, P): degP.;;;n},
del
p(j, Rn(z ;/, E, p == Rn(j, E, p) == inf{p(j, R): degR';;;n},
where the infima are over the set of polynomials P(z)
of degree degP:s;;;;n or rational functions R(z) of degree
deg R :s;;;; n, respectively (or over parts of these sets, dis-
tinguished by addition requirements). In fact, one deals
here with properties of solutions for a certain class of
extremal problems. In this context one may also con-
sider the study of other extremal problems on sets of
polynomials, rational functions and on certain classes
of analytic functions, as well as the study of analytic
properties of polynomials and rational functions (in
particular, obtaining inequalities between various norms
of these functions and their derivatives).
3) The study of the dependence of the rate of
decrease (to zero) of En(f, E, p) and Rn(f, E, p), as
n ~ 00, on the properties of f, E and p (the so-called
direct theorems of approximation theory) and the
dependence of the properties of f on the rate of
decrease of En(f, E, p) and Rn(f, E, p) to zero as n~oo
and on the properties of E and p (inverse theorems).
The study of approximation properties of well-known
methods in approximation theory (e.g., series in Faber
polynomials, various interpolation processes (cf. Inter-
polation process, as well as the search for new effec-
tive approximation methods are related to this direc-
tion.
4) The approximation of functions of several com-
plex variables. Here, basically, the same problems are
solved as in the case of one complex variable, but the
results and the methods for obtaining them differ, as a
rule, sharply from those used in the case of one vari-
able.
In the following some basic results have been listed.
I) The problem of the existence of a uniform approx-
imation by polynomials that is as good as one pleases
is solved by the Runge theorem (if f is analytic on E),
the Lavrent'ev theorem (if f is continuous on E), the
Keldysh theorem (if E is a closed domain, f is continu-
ous on E and analytic within E) and the Mergelyan
theorem (in the general case: E is a compact set, f is
continuous on E and analytic at interior points of E).
2) The problem whether it is possible to approximate
holomorphic fu~ctions on closed sets E in the extended
complex plane C is solved by Runge'S theorem. In the
study of approximating a function f in various spaces,
using the metric of these spaces, by rational functions,
an important role is played by characteristics of the set
e ::J C analogous to the analytic capacity y( e). In terms
of y(e) the problem of the description of all compact
sets E on which any continuous function can be
approximated with arbitrary accuracy by rational func-
231
APPROXIMATION OF FUNCTIONS OF A COMPLEX VARIABLE

tions is solved as follows: It is necessary and sufficient
that either
a) y(o(r, a)\E) = y(o(r, a = r
for any disc o(r, a)={z: I z-a I <r}, aEC,
that
b) lim y(o(r, a)\E) = 00
,--+0 r2
for any a EE (the equivalence of a) and b) expresses the
so-called 'instability' of the analytic capacity).
3) If E is bounded and Lebesgue-measurable and if
1~p <2, then the set of all rational functions is dense
in ~(E).
4) If P >0 and if G is a simply-connected domain
with a rectifiable Jordan boundary, then the family of
all polynomials in z is dense in the Smirnov class Ep(G)
if and only if G is a Smimov domain.
5) Let the complex-valued functions
I(Z),I/>I(Z), ... ,I/>n(z), n;;;'l, be continuous on a com-
pact set E c C. Among all generalized polynomials
P(z) = Ctq>t(z)+ ... +cnq>n(z)
(CI, ... ,Cn are arbitrary complex numbers) a general-
ized polynomial P o(z) deviates least from 1 in the
metric
Pe(j, P) = max{ I j(z)-P(z) I: zEE}
if and only if
min{Re[P(z)(Po(z)- j(z))]:
ZEE, Ij(z)-Po(z) I =Pe(j, Po)} ,;;; 0
for each P(z).
6) If E is a compact set with connected complement
G and if G has a Green function (for the first boundary
value problem for the Laplace equation) g(z, (0) with a
pole at 00, then for each Z E G and each polynomial
P(z) of degree n, the inequality
I P(z) I ,;;; M exp{ ng(z, oo)}
with M=max{ I P(z) I: Z EE} holds.
7) If E is a bounded non-degenerate continuum with
connected complement G and if I(z) is analytic at inte-
rior points of E and continuous on E with modulus of
8) In a number of cases for the approximation of
analytic functions various interpolation processes prove
effective, including Pade approximation and its generali-
zations.
r>O; or 9) For n ;;;'2, in en
there exist both non-closed Jor-
dan curves on which not every continuous function can
be uniformly approximated by polynomials in
(Zl, ... ,zn) to any degree of accuracy, and closed Jor-
dan curves on which polynomials uniformly approxi-
mate any continuous function. In CI this is impossible.
10) Up till now (1983) there are comparatively few
direct theorems on the approximation by rational func-
tions with free poles (i.e. without any condition on the
position of the poles of the approximating functions)
and a considerable amount of inverse theorems.
References
[I] GoNCHAROV, V.L.: The theory of interpolation and approxima-
tion offunctions, Moscow, 1954 (in Russian).
[2] WALSH, J.L.: Interpolation and approximation by rationalfunc-
tions in the complex domain, Arner. Math. Soc., 1969.
[3] SMIRNOV, V.I. and LEBEDEV, N.A.: Functions of a complex vari-
able: constructive theory, M.I.T., 1968 (translated from the Rus-
sian).
[4] DZYADYK, V.K.: Introduction to the theory of uniform approxi-
mation offunctions by polynomials, Moscow, 1977 (in Russian).
[5] RUSAK, V.N.: Rational functions as approximating tool, Minsk,
1979 (in Russian).
[6] GAMELIN, T.W.: Uniform algebras, Prentice-Hall, 1969.
[7] LEONT'EV, A.F.: Exponential series, Moscow, 1976 (in
Russian).
[8] Some questions in approximation theory, Moscow, 1963 (in Rus-
sian; translated from the English).
[9] KEWYSH, M.V.: Dokl. Akad. Nauk SSSR 4 (1936), 163-166.
[10] LAFRIENTIEFF, M. [M.A. LAVRENT'EV]: 'Zur Theorie der kon-
formen Abbildung', Trudy Mat. Inst. Steklov. 5 (1934), 159-
245.
[11] KOLMOGOROV, A.N.: 'Remark on the polynomials of P.L. Che-
byshev deviating least from a given function', Uspekhi Mat.
Nauk 3, no. 1 (1948),216-221 (in Russian).
[12] MERGELYAN, S.N.: 'Uniform approximation to functions of a
complex variable', Series and approximation 3 (1962), 294-391,
Arner. Math. Soc.. (Uspekhi Mat. Nauk 7, no. 2 (1952),31-
122)
[13] VlTUSHKIN, A.G.: 'The analytic capacity of sets in problems of
approximation theory', Russian Math. Surveys 22 (1967), 167-
200. (Uspekhi Mat. Nauk 22, no. 6 (1967), 141-199)
[14] DZHRBASYAN, M.M.: 'Some questions of the theory of
weighted polynomials in a complex domain', Mat. Sb. 36
(1955),353-440 (in Russian).

l l1~ JJ.
continuity w(~), then
[15] GONCHAR, A.A.: 'The rate of approximation of functions by
rational numbers and properties of numbers', in Proc. Internat.
En(j, E, pc) ,;;; C(j)w d n Congress Mathematicians Moscow, 1966, 1968, pp. 329-356 (in
Russian).
where [16] DOLZHENKO, E.P. and UL'YANOV, P.L.: Vestn. Moskov. Univ.
Ser. Mat. Mekh. 1 (1980), 3-13.
d(t) = max{min{ I ~-z I: ~EG, g(~, 00)=1n(1 +t)}: Z EaG}. [17J MERGELYAN, S.N.: 'On the approximation of functions of a
complex variable', in Mathematics in the USSR during 40
If the closed domain G is bounded by an analytic years: 1917-1957, Vol. I, Moscow, 1959, pp. 383-398 (in Rus-
curve T, then the condition sian).
[18] GONCHAR, A.A. and MERGELYAN, S.N.: History of paternal
En(j, E, pc) = O(n -p Q) mathematics, Vol. 1, Kiev, 1970, pp. 112-193 (in Russian).
is equivalent to the condition that f(P)(z) is Holder [19J TAMRAZOV, P.M.: Smoothness and polynomial approximation,
Kiev, 1975 (in Russian).
continuous of order <x, O<<x< 1, in G. The case when T [20J MEL'NIKOV, M.S. and SINANYAN, S.O.: Contemporary problems
is a piecewise-smooth curve with corners has been stu- in mathematics, Vol. 4, Moscow, 1975, pp. 143-250 (in Rus-
died. sian).
E. P. Dolzhenko

232
 APPROXIMATION OF FUNCTIONS OF SEVERAL REAL VARIABLES
Editorial comments. Let EC C be a compact set. Let
A(E) denote the set of functions f that are continuous on E
and analytic at interior points (if any) of E. Denote by peE)
(respectively, R(E) the set of functions fEA(E) that can be
uniformly approximated on E by polynomials (respectively,
rational functions). Mergelyan's theorem states: 1)
A(E)=P(E) if and only if CE, the complement of E, is con-
nected; 2) A(E) = R(E) if CE is finitely-connected. Examples
of compact sets E for which A(E)=I=R(E) are known (e.g.
the Schweizer Kase of A. Roth, cf. [A 1]). The problem
arises of characterizing those E for which A(E) = R(E). The
solution was given by E. Bishop and A.G. Vitushkin
(independently) for compact sets without interior pOints, and
by Vitushkin for general compact E. Vitushkin's theorem can
be found in [A1].
The sets A(E), peE) and R(E) are uniform algebras, e.g.
function algebras endowed with the sup-norm (cf. Algebra
of functions; Uniform algebra). For studies stressing this
aspect see [6] and [M].
Another direction of research is to study not approxima-
tion by linear combinations of 1, Z, Z2, . . . , or
... ,Z-1, 1, z, .. (polynomials or rational functions) but by
linear combinations of powers ZA" or of exponentials eiA,Z,
where {An} cR, on sets ECC, mostly on curves satisfying
some 'oscillation condition'. Results of this type related to
the Muntz theorem; Lacunary power series; the
Paley-Wiener theorem, zero sets of analytic functions,
etc. (cf. [AS]).
In cn approximation problems depend essentially on the
geometry of the domain under consideration (this is related
to the so-called Levi problem of characterizing domains of
holomorphy). One way to approach approximation problems
a
is via H6rmander's mechanism (cf. [A6]). An example of
an approximation theorem is Kerzman's theorem: Let nc C2
be a strongly pseudo-convex domain with sufficiently
smooth boundary (C 5 suffices). Then every function fez)
holomorphic on n and continuous on Q can be uniformly
approximated on Q by functions f} that are holomorphic ~
A
some (strongly pseudo-convex) domain n containing n,
[A7].
References
[A 1] GAIER, D.: Vorlesungen iiber Approximation im Komplexen,
Birkhauser, 1980.
[A2] BUCK, R.c.: 'Survey of recent Russian literature on approxi-
mation', in R.E. Langer (ed.): On numerical approximation,
Univ. of Wisconsin Press, 1959, pp. 341-359.
[A3] KOREVAAR, J.: 'Polynomial and rational approximation in the
complex domain', in J.G. Clunie (ed.): Aspects of contem-
porary complex analysis, Acad. Press, 1980, pp. 251-292.
[A4] STOUT, E.L.: The theory of uniform algebras, Bogden &
Quigley, 1971.
[A5] REDHEFFER, R.M.: 'Completeness of sets of complex
exponentials', Adv. in Math. 24 (1977), 1-62.
[A6] HORMANDER, L.: An introduction to complex analysis in
several variables, North-HOiland, 1973.
[A7] KERZMAN, N.: 'Holder and LP estimates for solutions of
au= f on strongly pseudo-convex domains', Commun. Pure
Appl. Math. 24 (1971), 301-380.
AMS 1980 Subject Classification: 30E10, 41 A25
APPROXIMATION OF FUNCflONS OF SEVERAL
REAL VARIABLES - The case when the function f to
be approximated depends on two or more variables:
f(t) = f(t], ... ,tm ), m ;;;'2.
(See Approximation of functions.) In comparison with
the one-dimensional case, the study of problems on the
approximation of functions of m (m ;;;;'2) variables is
highly complicated because of essentially new cir-
cumstances that arise related to dimension. First of all
this applies to the domain on which the approximation
is carried out. Simply-connected compact sets (in the
one-dimensional case intervals) in Rm (even in the
plane) can have a large variety of shapes; therefore it
becomes necessary to classify the domains according to,
for instance, the smoothness properties of their boun-
daries. Complications also arise in the description of
difference-differential properties of functions of m vari-
ables. Generally speaking, these properties can be dif-
ferent in different directions, and their characterization
ought to take into account both the geometry of the
domain on which the function is defined and the
behaviour of the function while approaching the boun-
dary, so that the study of boundary properties of the
function is of great importance. If one attempts to sim-
plify the solution of an approximation problem by
passing to a domain with a simpler structure, the prob-
lem that arises is to extend a function f from a domain
Q cRm into a canonical domain Ql (for instance a
parallelepiped or the whole space Rm) containing Q
while preserving certain smoothness properties (see
Extension theorems). This kind of problems is closely
connected with imbedding theorems, and also with
questions arising in the solution of the boundary value
problems of mathematical physics.
An increase in the number of independent variables
complicates, of course, also the tools of approximation,
because when the dimension is increased, with it, e.g.,
the degrees of the polynomials increase. An algebraic
polynomial of degrees n I, . . . ,nm in the variables
t I, . . . ,tm , respectively, has the form
Pn" ... ,nm(t I , " " - ~ak"
tm ) -.c... l' .. lmm,
... ,km 1 (1)
where ak" ... ,km are real coefficients. The summation is
carried out with respect to k P = I, . . . ,m, from 0 to
pO
n so that, e.g., the subspace of polynomials of degree
pO
3 in each of the m variables has dimension 4m. Some-
times the overall or total degree n of the polynomial is
fixed; then the summation in (1) is carried out over all
indices satisfying the inequality O';;;;k 1 + ... +km.;;;;n.
A real trigonometric polynomial of degrees n I, . . . ,nm
in the variables t I, . . . ,tm can be written in the form
Tn" ... ,nJt l , .. , ,tm) = ~ak, ..... k,exp [(~kptp],
p=1
where complex coefficients ak, ... ,k, with indices of
opposite sign are complex conjugate, and the summa-
tion is carried out with respect to k P = I, ... ,m, pO
233
APPROXIMATION OF FUNCTIONS OF SEVERAL REAL VARIABLES

from - n" to n". Such a polynomial can be Let Q be an arbitrary open set in Rm (in particular,
represented in the form of a linear combination of all Q= Rm), let e be a unit vector in Rm, let h >0 and let
possible products of the form <Pk,(tl) ... <Pk.(tm), where Qhe be the set of points t E Q such that the segment
<Pk.(t,,) is either sink"t" (OE;;k,,=E;;;n,,) or cosk"t" [t, t+he]EQ. If fELp(Q) and l=e;;;p=E;;;oo, then the quan-
(O=E;;;k,,=E;;;n,,). Multi-dimensional splines have a wide tity
field of applications; they consist of 'pieces' of alge- weif; 8)~(Q) = ~~ II f(t +he)- f(t) II ~(Q.,.)
braic polynomials in m variables 'tied' together in order
is called the modulus of continuity of the function
to obtain specified smoothness properties. In the case
f(t I, . . . ,tm) in the direction e with respect to the metric
m = 2, splines of the simplest form are formed by poly-
of Lp(Q). The quantity
nomials 'tied' together along straight lines parallel to
the coordinate axes. As a tool of approximation func- wif; 8)~(Q) = sup weif; 8)~(Q)
e
tions g(tI. ... ,tm ) are used which are polynomials or
is called the modulus of continuity of f in Lp(Q).
splines only in one of the variables. For the approxi- _ In the periodic case, the best _approximation
mation of non-periodic functions given on the whole
En ..... ,n.(fi,(Rm of a function fELp(Rm) having
space Rm (or on an unbounded subset of Rm) one may
(Sobolev generalized) partial derivatives
use entire functions of exponential type. These can be
represented in the form of sums of absolutely conver- D~f = L f E i,,(Rm)
gent power series at:
Gn.. ... .nJt;, ... ,tm) = ~ ak .. ... k. tt, ... t~ (2)
(where the r,,;;;;'O are integers and DOf=f,
k."O. p= 1, ... ,m) by trigonometric polynomials Tn" ... ,n.,
p=l ... m
satisfies the inequality
under the condition that for any >0 and all complex _ m
tI. ... ,tm one has E....... n.if)"i,(R") .;;;;; M~n;rwe.(D:f;n;I)i,(R.), (3)
m .=1
I Gn .......Jt\, ... ,tm) I .;;;; M,exp ~(n.+()1 t.l,
where e" is the unit vector directed along t", and where
where the constant M. depends only on (see [1]). It the constant}1 depends neither on f nor on n". For a
should be noted that, unlike polynomials, the function function fELp(Rm) having generalized partial deriva-
(2) is determined by an infinite number of parameters. tives
In the multi-dimensional case one has Weierstrass'
theorem on the possibility of approximating a function
fEC(Q), i.e. continuous on_abounded and closed set of order r=rl + ... +rm (where r=(rI. ... ,rm, the
Q CRm, or a function fEC(Rm), continuous on the
following inequalities hold:
whole space Rm with period 2'1T in each variable, with E..... n(f)i,(R) .;;;;;~ ~ w(Drf;n-I)i,(R)' (4)
arbitrary accuracy by means of algebraic (respectively, n'I+"'+,,,,;::::7
trigonometric) polynomi~s. A similar result holds in If
w(Drf; 8)i.,(Rm ) Kif', 0<0:< 1,
the spaces Lp(Q) and Lp(Rm) (in the periodic case) .;;;;;

(1 =E;;;p < 00). General results and theorems concerning i.e. if the function Drfsatisfies a Holder condition, then
properties of the best approximation, the existence and
ff)-.";;:' ~
uniqueness, the characteristic properties of functions of
E- n... nV ~(R) '"'" r+a' r -0
- , 1, ... ., 0< a..."
.,;;:,1 . (5)
n
best approximation, and general relations of duality In the last case the inverse theorem asserts that if for a
when approximating by (means of) a convex set (of function fELp(Rm) inequality (5) hol,?s for all
functions) and, in particular, by a subspace, can be n = 1,2, ... , then the derivatives DrfELp(Rm) exist
extended to normed linear spaces of functions of m and satisfy for any bERm the inequalities
variables (see [3] and [4]). However, to obtain explicit
formulations of these theorems, taking into account a II Drf(t+h)-Drf(t) 1Ii.,(R m) .;;;;; K I h la (6)
specific metric and specific properties of the approxi- for 0<0:< 1; whereas
mating subspace in the multi-dimensional case, involves II Drf(t+h)-2Drf(t)-Drf(t-h) 1Ii.,(Rm
) .; ; ; K Ih I (7)
great difficulties.
Quite thorough investigations were carried out on for 0: = 1. Here K does not depend on the length
questions concerning the connection between the I b I =(hr + ... +hm )I/2 of the vector
smoothness properties of a function of sev.eral variables b=(hl' ... ,hm ).
and the rate of decrease of its best approximation by Theorems analogous to those given above also hold
algebraic and trigonometric polynomials, as well as by for non-periodic functions fELp(Rm) if entire functions
entire functions. of exponential type are used as tools of approximation.

234

The above-mentioned results can be extended also to
classes of functions whose smoothness is described in
terms of moduli of continuity (moduli of smoothness)
of a higher order (see [1]).
In the case of approximating functions /ELP(Q) by
algebraic polynomials P n" ... ,n~ on a bounded paral-
lelepiped (as well on some other bounded sets) direct
theorems similar to (3), (4) and (5) have been proved.
Inverses of these theorems, like in the case of functions
defined on a bounded interval, are possible only on a
proper subset Q 1 of Q. Inverse theorems are known
assuming a better order of approximation in the neigh-
bourhood of boundary of Q (see [13]); there are also
direct theorems asserting the possibility of achieving
such an improvement in a neighbourhood of angular
points (see [14]). Necessary and sufficient conditions
for a function / to belong to a class H'c+a(Q) (defined
in the metric of C(Q)) by conditions similar to the con-
ditions (6) and (7)) at the cost of raising the order of
approximation in a neighbourhood of the boundary
(like in the one-dimensional case) are unknown (1983).
However, one has the following results of a negative
character (see [13]). Let Q={t: tER2 , I t I ~1}. There
does not exist a sequence A~a)( I t I), n = 1, 2, ... ;
0<0:< 1, of functions in Q having the two following
properties:
1) For any function /EH'C(Q) there are a constant
M and a sequence of polynomials
n=l,2, ... ,
such that
I I(t)- Pn(t) I .;;; MA~")( I t I), tEQ.
2) The existence of a constant M>O and of a
sequence of polynomials Pn(t) satisfying (8) implies that
/EH'C(Q) for any function/defined on Q.
To illustrate the specific nature of approximation of
functions of several variables, one can mention the fol-
lowing ~esults.
Let En"n, (j)x be the best approximation of a 2'lT-
periodic function / of two variables by trigonometric
I:0ly~omials Tn"n, in the metric of X (X=C(R2) or
X=Lp (R 2 )), _and let En"oo(j)X be the best approxima-
tion of / in X by functions Tn" 00' i.e. by trigonometric
polynomials
.
of degree at most n 1 in the variable I 1 and
~th coefficients that are functions of (2' The quantity
E oo,n, (j)x is defined similarly.
If I <p < 00, one has the inequalities
En,.n, (f)L,(R') .;;; Ap {En,.oo(f)L,(R') + E oo.n, (f)L,(R')},
wher~ Ap _depends only on p.
If X=L 1(R 2) or X=C(R 2), then
APPROXIMATION OF FUNCTIONS OF SEVERAL REAL VARIABLES
.;;; A In(2+min{nl, n2})(En"oo(f)x+Eoo,n,(f)x),
where A is an absolute constant, and the factor
1n(2 + min {n 1 , n 2 }) in (9) cannot be replaced by
another factor tending slower to infinity as
min{nl' n2}~00 (see [15]).
Fundamental peculiarities arise in problems of inter-
polation for functions of several variables. For instance,
contrary to the one-dimensional case, the existence of
an algebraic interpolation polynomial depends substan-
tially on the position of the interpolation nodes.
Nevertheless, effective methods of constructing polyno-
mials and splines interpolating a function /(/1, ... ,1m)
on a distribution of nodes chosen in a definite way,
have been established (see Interpolation). For multi-
dimensional interpolation by splines, in a number of
cases order estimates of the approximation errors, for
the function / as well as for its partial derivatives, have
been found; in this direction, two-dimensional splines
of low degree, as well as local (Hermitian) splines of
arbitrary degree, have been studied in more detail (see
[7], [10] - [12]). Among the other linear methods for
approximating functions of several variables, multiple
Fourier sums and their various means have been stu-
died most extensively. Here, order estimates of the
approximation with respect to classes of functions are
known, and in some case asymptotically exact results
were obtained (see [5], [6] and [8]).
References
[I] NIKOL'sKIi, S.M.: Approximation offunctions of seJieral variables
and imbedding theorems, Springer, 1975 (translated from the
Russian).
(8) [2] GUTER, R.S., KUDRYAVTSEV, L.D. and LEVITAN, B.M.: Ele-
ments of the theory offunctions, Moscow, 1963, 106-198 (in
Russian).
[3] KORNEicHUK, N.P.: Extremal problems in approximation theory,
Moscow, 1976 (in Russian).
[4] TIKHOMIROV, V.M.: Some questions in approximation theory,
Moscow, 1976 (in Russian).
[5] DZYADYK, V.K.: Introduction to the theory of uniform approxi-
mation offunctions by polynomials, Moscow, 1977 (in Russian).
[6] TIMAN, A.F.: Theory of approximation offunctions of a real
variable, Pergamon, 1963 (translated from the Russian).
[7] STECHKlN, S.B. and SUBBOTIN, YU.N.: Splines in numerical
mathematics, Moscow, 1976 (in Russian).
[8] STEPANETS, A.I.: Uniform approximation by trigonometric
polynomials. Linear methods, Kiev, 1980 (in Russian).
[9] LAURENT, P.l.: Approximation et optimisation, Hermann, 1972.
[10] AHLBERG, l.H., NILSON, E.N. and WALSH, l.F.: Theory of
splines and their applications, Acad. Press, 1967.
[II] ZAV'YALOV, Yu.S., KVASOV, B.1. and MIROSHICHENKO V L .
Methods of spline functions, Moscow, 1980 (in Russian). . ..
[12] VARGA, R.S.: Functional analySiS and approximation theory in
numerical analysis, Reg. Conf. Series in App!. Math., 3, SIAM,
1971.
[13] NIKOL'sKIi, S.M.: 'Approximation of functions of several vari-
ables by polynomials', Siberian Math. 1. 10, no. 4 (1969),792-
799. (Sibirsk. Mat. Zh. 10 (1969), 1075-1083)
[14] BRUDNYi, YU.A.: 'Approximation of functions defined in a
convex polyhedron', Soviet Math. Doklady 11, no. 6 (1970),
(9) 1587-1590. (Dokl. Akad Nauk SSSR 195 (1970), 1007-1009)
235
APPROXIMATION OF FUNCTIONS OF SEVERAL REAL VARIABLES
[15] TEML'YAKov, V.N.: 'Best approximations for functions of two
variables', Soviet Math. Doklady 16, no. 4 (1975), 1051-1055.
(Dokl. Akad. Nauk SSSR 223 (1975), 1079-1082)
V. N. Konovalov
N.P. Korneichuk
References
[A1] CHENEY, E.W.: 'Four lectures on multivariate approximation',
in S.P. Singh, J.HW. Burry and B. Watson (eds.): Approxi-
mation theory and spline functions, Reidel, 1984, pp. 65-87.
[A2] SCHEMPP, W. and ZELLER, K. (EDS.): Multivariate approxima-
tion theory, Birkhauser, 1979.
[A3] SCHEMPP, W. and ZELLER, K. (EDS.): Multivariate approxima-
tion theory /I, Birkhauser, 1982.
[A4] SCHEMPP, W. and ZELLER, K. (EDS.): Multivariate approxima-
tion theory /II, Birkhauser, 1985.
[A5] ScHEMPP, W. and ZELLER, K. (EDS.): Constructive theory of
functions of several variables, Springer, 1977.
[A6] LORENTZ, G.G.: Approximation of functions, Holt, Rinehart
and Winston, 1966.
[A7] GOLOMB, M.: 'Approximation by functions of fewer variables'
in R.E. Langer (ed.): On numerical approximation, Univ. of
Wisconsin Press, 1959, pp. 275-327.
[A8] SCHUMAKER, L.L.: Spline functions. Basic theory, Wiley,
1981.
AMS 19S0 Subject Classification: 41-XX, 41 A03,
41A25,41A27
APPROXIMATION ORDER - The order of the error
of approximation as a variable quantity, depending on
a continuous or discrete argument 'T, relative to another
variable 1f>('T) whose behaviour, as a rule, is assumed to
be known. In general, 'T is a parameter that is a numeri-
cal characteristic of the approximating set (e.g. its
dimension) or of the method of approximation (e.g. the
interpolation step). The set of values of 'T may, more-
over, have an infinite or finite limit point. The function
1f>('T) is most often a power, an exponential or a loga-
rithmic function. The modulus of continuity (cf. Con-
tinuity, modulus of) of the approximated function (or
that of some derivative of it) or a majorant of it may
figure as If>('T).
The approximation order is characterized both by the
properties of the approximation method, as well as by a
definite property of the approximated object, e.g. the
differential-difference properties of the approximated
function (cf. Approximation of functions, direct and
inverse theorems).
In numerical analysis, the approximation order of a
numerical method having error O(h m ), where h is the
step of the method, is the exponent m.
References
[I] GONCHAROV, V.L.: The theory of interpolation and approxima
tion of functions, Moscow, 1954 (in Russian).
[2] TIMAN, A.F.: Theory of approximation offunctions of a real
variable, Pergamon, 1963 (translated from the Russian).
[3] BAKHVALOV, N.S.: Numerical methods: analysis,. algebra, ordi
nary differential equations. Mir. 1977 (translated from the Rus-
sian).
N.P. Korneichuk
v.P. Motornyi
236
Editorial comments.
References
[A1] MOllER, M.W.: Approximationstheorie, Akad. Ver-
lagsgesellschaft, Wiesbaden, 1978.
[A2] LORENTZ, G.G.: Approximation of functions, Holt, Rinehart
and Winston, 1966.
AMS 1980 Subject Classification: 41 A25
APPROXIMATION THEORY - The branch of
mathematical analysis studying methods for approxi-
mating some mathematical objects by others and also
studying questions related to the research and estima-
tion of the error that arises here.
The main contents of approximation theory concerns
the approximation of functions. Its foundations are laid
by the work of P.L. Chebyshev (1854 - 1859) on best
uniform approximation of functions by polynomials
and by K. Weierstrass, who in 1885 established that in
principle it is possible to approximate a continuous
function on a finite interval by polynomials with arbi-
trary pre-given error. The development of approxima-
tion theory was to a large extent determined by the
fundamental work of H. Lebesgue, Ch.J. de la Vallee-
Poussin, S.N. Bernstein [S.N. Bernshtein], D. Jackson,
J. Favard, A.N. Kolmogorov and S.M. Nikol'skii on
the approximation of functions and function classes.
With the development of functional analysis, many
problems in approximation theory were considered in
the most general setting, e.g. as the approximation of
elements of an arbitrary linear normed space X. Three
classes of problems arose here, corresponding more or
less to the main chronological stages of development of
research in approximation theory.
1. The approximation of a fixed element x E X by ele-
ments of a fixed set 9C c X. If
E(x, 91) = inf II x-u I
UE9(
is taken as the approximation measure, i.e. the best
approximation of x by 91, then, along with the investi-
gation and estimation of E(x, 9C), questions on the
existence of an element of best approximation Uo E9C
(for which I x - Uo I = E(x, 9C, its uniqueness and
characteristic features arise. Any operator A mapping X
into 9C gives rise to an approximation method with
error I x - Ax II. If 9C is a linear manifold, linear
operators are of particular importance. For sequences
{An} of such operators, the question of the conditions
of convergence Anx--n for any x EX arises.
2. The approximation of a fixed set we C X by ele-
ments of another fixed set 91 C X. The best approxima-
tion here is expressed by
E(IJJI,91) = .suo
XE::!1,
E(x, 91),
which glves the minimal possible error estimate when

approximating an arbitrary x EX by elements from 91.
In concrete cases, the problem consists of estimating, or
expressing exactly, E('iJR,91) by characteristic properties
of the given sets 'iJR and 91. If the approximation is
established by an operator A, the supremum
~~~ II x-Ax II
is investigated, as well as (if 91 is a linear manifold)
6'(IJR,9C) inf sup I x -
= AXe91XE~ Ax II,
where the infimum is taken over all linear operators A
mapping X into 91. A linear operator realizing this
infimum (if it exists) gives rise to a best linear method
of approximation. The case lS'('iJR, 91)=E('iJR, 91) is of
particular interest.
3. Best approximation of a fixed set 'iJR C X by a
given class {91} of approximating sets in X. It is
assumed that in {91} there are, in a definite sense,
'equally-expensive' classes, e.g. containing the same
amount of elements or having the same dimension. The
first case leads to the (-entropy of 'iJR (relative to X),
the second - to the problem of calculating the width of
'iJR (in X), in particular,
dN(IJR, X) = inf
'liN
(IJR, 9C N ), (1)
and
d~(IJR, X) = inf 6' (IJR, 9C N ), (2)
9lN
where the infimum is taken over all subspaces 91N in X
of fixed dimension N (or over all possible translations
91N +a of it). Thus in (1), (2) the problem is to deter-
mine the best (respectively the best linear) approxima-
tion tool of dimension N for 'iJR.
References
[I] ACHIEZER, N.!. [N.!. AxHIEZER]: Theory of approximation, F.
Ungar, 1956 (translated from the Russian).
[2] GONCHAROV, V.L.: The theory of interpolation and approxima-
tion offunctions, Moscow, 1954 (in Russian).
[3] DZYADYK, V.K.: Introduction to the theory of uniform approxi-
mation offunctions by polynomials, Moscow, 1977 (in Russian).
[4] NIKOL'SKIi, S.M.: Approximation offunctions of several variables
and imbedding theorems, Springer, 1975 (translated from the
Russian).
[5] KORNEicHUK, N.P.: Extremal problems in approximation theory,
Moscow, 1976 (in Russian).
[6] TIKHOMIROV, V.M.: Some questions in approximation theory,
Moscow, 1976 (in Russian).
[7] TIMAN, A.F.: Theory of approximation offunctions of a real
variable, Pergamon, 1963 (translated from the Russian).
[8] RICE, J.R.: The approximation offunctions, 1-2, Addison-
Wesley, 1964-1968.
[9] DAVIS, PJ.: Interpolation and approximation, Blaisdell, 1965.
[10] LORENTZ, G.G.: Approximation offunctions, Holt, Rinehart
and Winston, 1966.
N.P. Kornefchuk
VP. Motornyi
Editorial comments. There are many good books on
approximation (and interpolation) theory. Below some of
them are mentioned. [A1] contains extensive notes (related
ARBITRA nON SCHEME
to the history of theorems and methods), as well as a useful
bibliography. The more recent book [A2] also contains an
extensive bi bl iog raphy .
References
[A1] CHENEY, E.W.: Introduction to approximation theory, Chel-
sea, reprint, 1982.
[A2] HOLLAND, A.S.B. and SAHNEY, B.N.: The general problem of
approximation and spline functions, R.E. Krieger, 1979.
[A3] NATANSON, !'P.: Constructive theory of functions, F. Ungar,
1964-1965 (translated from the Russian).
[A4] SINGER, I.: Best approximation in normed linear spaces by
elements of linear subspaces, Springer, 1970 (translated
from the Rumanian).
[AS] PINKUS, A.: n-widths in approximation theory, Springer, 1985
(translated from the Russian).
AMS 1980 Subject Classification: 41 A50, 41 A46,
41A65
ARABIC NUMERALS - The traditional name given
to the mathematical symbols: 0, 1, 2, 3, 4, 5, 6, 7, 8, 9.
Using these symbols, every number can be expressed in
the decimal system (d. Numbers, representations of).
These numerals arose in India (not later than the fifth
century), in Europe they became known in the lO-th to
13-th centuries through Arabic works (hence the name).
BSE-3
AMS 1980 Subject Classification: 01-XX
ARBITRATION SCHEME - A rule by which each
arbitration game (d. Cooperative game) is put into
correspondence with a unique outcome of the game is
called an arbitration solution. The first arbitration
scheme was considered by J. Nash [1] for the case of a
two-person game. Let R = { u = (u \, ... ,un)} be the set
of outcomes, let d = (d \, ... ,dn ) be the status quo
point, i.e. the point corresponding to the situation in
which no cooperative outcome is realized, let [R, d] be
the corresponding arbitration game and let u be an
arbitration solution of it. An outcome u is called a
Nash solution if
II(u;-d;) = maxII(u;-d;),
i UER i
Only a Nash solution satisfies the following axioms: 1)
if f is a linear non-decreasing mapping then fo is an
arbitration solution of the game [jR, fd] (invariance
with respect to utility transformations); 2) u~d, u ER
and there is no u ER such that u;:. u (Pareto optimal-
u
ity); 3) if R' CR, d' =d, u ER', then =u (indepen-
dence of irrelevant alternatives); and 4) if d; =dj ,
i, j = 1, ... ,n, and R is symmetric, then u; = Uj'
i,j=l, ... ,n (symmetry).
Another arbitration scheme for an n-person game
with characteristic function v and player set
N = {1, ... ,n} was given by L.S. Shapley [2]. The
Shapley solution cj>( v) = (q,\ (v), ... ,q,n (v)), where
q,;(v) = ~ Yn(s)[v(S)-v(S \ {i})],
seN
237
ARBITRA TION SCHEME
'Yn(S)=(S -l)!(n -S)! / n! and S is the number of ele-
ments of the set S, also satisfies the axiom of sym-
metry, but, moreover, ~.(Mv)=v(N), and for any two
I
games u and v the equality <p(u +v)=<p(u)+<p(v) holds.
Arbitration schemes with interpersonal utility comparis-
ons have also been considered [3].
The arbitration schemes of Nash and Shapley were
generalized by J.e. Harsanyi [4]. A Harsanyi solution
satisfies, apart from the four axioms of Nash, the two
axioms: I) the solution depends monotonically on the
initial demands of the players; and 2) if u' and u" are
two solutions, then Ii, defined by
u ~ min(u;,
ieN
u;*),
is also a solution if and only if Ii belongs to the boun-
dary of the set R.
Under suitable conditions an arbitration scheme
depends continuously on the parameters of the game.
References
[l] NASH, J.: 'The bargaining problem', Econometrica 18, no. 2
(1950), 155-162.
[2] SHAPLEY, L.S.: 'A value for n-person games', in H.W. Kuhn,
A.W. Tucker and M. Dresher (eds.): Contributions to the theory
of games, Vol. 2, Princeton Univ. Press, 1953, pp. 307-317.
[3] RAiFFA, H.: 'Arbitration schemes for generalized two-person
games', in H.W. Kuhn, A.W. Tucker and M. Dresher (eds.):
Contributions to the theory of games, Vol. 2, Princeton Univ.
Press, 1953, pp. 361-387.
[4] HARsANYI, J.e.: 'A bargaining model for the cooperative n-
person game', in H.W. Kuhn, A.W. Tucker and M. Dresher
(eds.): Contributions to the theory of games, Vol. 4, Princeton
Univ. Press, 1959, pp. 325-355.
E.I. Vilkas
Editorial comments. Arbitration schemes are also called
bargaining schemes and a Nash solution is also called a
bargaining solution. The Shapley solution vector <p is also
called the Shapley value. For other, more recent, bargain-
ing schemes, such as the Kalai- Smorodinsky solution and
Szidarovsky's generalization of the concept of a Nash solu-
tion, the reader is referred to [A 1], [A2] , respectively [A6].
For further developments concerning Harsanyi solutions, cf.
[A3]. Some authors distinguish between bargaining
schemes and arbitration schemes. Then the Nash scheme is
a bargaining scheme and the Shapley one an arbitration
scheme, [A5].
References
[A 1] KALAl, E. and SMORODINSKY, M.: 'Other solutions to Nash's
bargaining problems', Econometrica 43 (1975), 513-518.
[A2] ROTH, A.E.: Axiomatic models of bargaining, Lect. Notes in
Econom. and Math. Systems, 170, Springer, 1979.
[A3] HARSANYI, J.e.: Papers in game theory, Reidel, 1982.
[A4] AUMANN, R.I. and SHAPLEY, L.S.: Values of non-atomic
games, Princeton Univ. Press, 1974.
[A5] RApOPORT, A.: N-person game theory: Concepts and appli-
cations, Univ. of Michigan Press, 1970, p. 168.
[A6] SZEP, 1. and FORG6, F.: Introduction to the theory of games,
Reidel, 1985.
[A7] VOROB'EV, N.N.: Game theory, Springer, 1977 (translated
from the Russian).
AMS 1980 Subject Classification: 90012
ARc, simple arc, Jordan arc - A part of a curve
comprised between two of its points (and containing no
multiple points). An arc in the plane is often defined
by specifying the coordinates of its points as continu-
ous functions x = <p(t), y =#..1) of some parameter t,
a";;;;'t ";;;;'b; it is assumed in so doing that different values
of t correspond to different points.
Editorial comments. In general any homeomorphic
image of the unit interval is called an arc, not to be con-
fused with a path: A continuous function with as domain the
unit interval.
A related concept is that of a Hausdorff arc (or sometimes
also arc): A continuum (a compact and connected topologi-
cal space) that is irreducibly connected between two of its
pOints (any closed and connected subset containing these
two points must be equal to the whole space).
References
[A1] KURATOWSKI, K.: Topology, 2, Acad. Press, 1968 (translated
from the French).
AMS 1980 Subject Classification: 30A99, 51 N99,
53A04, 54FXX
ARC, CONTACfLESS (FREE) - A smooth curve
without self-intersections in the phase plane of a two-
dimensional autonomous system of differential equa-
tions
x= P(x,y), y= Q(x,y), (*)
at each point of which the phase velocity vector of the
system (cf. Phase velocity vector) is defined, is non-
zero and is not a vector tangent to the curve. The con-
cept was introduced by H. Poincare [1), and is exten-
sively employed in the qualitative theory of differential
equations (2). Thus, it is possible to draw a contact less
segment of sufficiently small length through an arbi-
trary ordinary point of a trajectory of the system (*). A
con tactless arc is characterized by the fact that all the
trajectories of system (*) intersecting the curve intersect
it in the same direction. If the derivative along the flow
of the system (*) (cf. Differentiation along the flow of a
dynamical system) at each point of the smooth curve
does not vanish, this curve is a contactless arc. A closed
contactless arc is said to be a contactless cycle.
References
[I] POINCARE, H.: 'Memoire sur les courbes definies par une
equation differentielle I - IV', in Oeuvres de H. Poincare, Vol.
I, Gauthier-Villars, 1916, pp. 3-222.
[2] LEFSCHFTZ, S.: Differential equations: geometric theory, Intersci
ence, 1962.
N.Kh. Rozov
AMS 1980 Subject Classification: 34C35, 58F99
ARC FUNCTION, arcus junction - A function inverse
to a trigonometric function, i.e. one of the functions:
arcsine, arccosine, arctangent, arccotangent, arc secant,
arccosecant. Cf. Inverse trigonometric functions.

Editorial comments.
References
[1] ABRAMOWITZ, M. and STEGUN, J.A.: Handbook of mathemati-
cal functions, National Bureau of Standards, 1964. Dover
reprint, 1972.
AMS 1980 Subject Classification: 26A09
ARCHIMEDEAN AXIOM - An axiom, originally for-
mulated for segments, which states that if the smaller
one of two given segments is marked off a sufficient
number of times, it will always produce a segment
larger than the larger one of the original two segments.
This axiom can be formulated in an analogous manner
for surfaces, volumes, positive numbers, etc. In general,
the Archimedean axiom applies to a given quantity if
for any two values A and B of this quantity such that
A <B it is always possible to find an integer m such
that Am> B. The axiom forms the base of the process
of successive division in arithmetic and in geometry (cf.
Euclidean algorithm). The importance of the
Archimedean axiom only became fully apparent after
the discovery, in the 19-th century, of magnitudes to
which it does not apply - so-called non-Archimedean
quantities (cf. Quantity; see also Archimedean group;
Archimedean ring; Archimedean class).
The axiom was explicitly formulated by Archimedes
(third century B.C.) in his work The sphere and the
cylinder; it had been previously utilized by Eudoxus of
Cnides, and is for this reason also called Eudoxus'
axiom.
BSE-3
AMS 1980 Subject Classification: 01 A20, 06F05
ARCHIMEDEAN BODIES - A synonym for Semi-
regular polyhedra.
AMS 1980 Subject Classification: 52A25, 51 M20
ARCHIMEDEAN CLASS - A class resulting from the
subdivision induced by the Archimedean equivalence
relation on a totally ordered semi-group. This
equivalence is defined as follows. Two elements a and b
of a semi-group S are called Archimedean equivalent if
one of the following four relations is satisfied:
a"';;;;b"';;;;a n , b"';;;;a",;;;;b n,
an",;;;;b"';;;;a, bn"';;;;a"';;;;b;
which amounts to saying that a and b generate the
same convex sub-semi-group in S. Thus, the subdivi-
sion into Archimedean classes is a subdivision into
pairwise non-intersecting convex sub-semi-groups.
Moreover, each subdivision into pairwise non-
intersecting convex sub-semi-groups, can be extended
to a subdivision into Archimedean classes.
The Archimedean equivalence on a totally ordered
group is induced by the Archimedean equivalence of its
ARCHIMEDEAN SEMI-GROUP
positive cone: It is considered that a,..."b if there exist
positive integers m and n such that
I a I < I b Im and I b I < I a In,
where
Ixl = max{x, X-I}.
The positive cone of an Archimedean group consists of
a single Archimedean class.
O.A. Ivanova
AMS 1980 Subject Classification: 06F05
ARCIDMEDEAN GROUP - A partially ordered group
within which the Archimedean axiom is valid: If an <b
for all integers n (a and b are elements of the
Archimedean group), then a is the unit of the group (in
additive notation: If na <b for all integers n, then
a =0). Totally ordered Archimedean groups may be
described as follows (Holder's theorem): A totally
ordered group is Archimedean if and only if it is iso-
morphic to some subgroup of the additive group of real
numbers with the natural order. Thus, the additive
group of all real numbers is in a certain sense the larg-
est totally ordered Archimedean group. All lattice
ordered Archimedean groups are commutative. Totally
ordered groups without non-trivial convex subgroups
(cf. Convex subgroup) are Archimedean.
References
[1] BIRKHOFF, G.: Lattice theory" Amer. Math. Soc., 1967.
[2] KOKORIN, A.1. and KoPYTov, A.M.: Fully ordered groups,
Israel Progr. Sci. Trans!., 1974.
[3] fuCHS, L.: Partially ordered algebraic systems, Pergamon, 1963.
A.!, Kokorin
V.M. Kopytov
AMS 1980 Subject Classification: 06F15
ARCIDMEDEAN RING - A partially ordered ring the
additive group of which is an Archimedean group with
respect to the given order. An Archimedean totally
ordered ring R is either a ring with zero multiplication
(i.e. xy =0 for all x and y in R) over an additive group
which is isomorphic to some subgroup of the group of
real numbers, or else is isomorphic to a unique subring
of the field of real numbers, taken with the usual order.
An Archimedean totally ordered ring is always associa-
tive and commutative.
References
[I) FUCHS, L.: Partially ordered algebraic systems, Pergamon Press,
1963.
o.A. Ivanova
AMS 1980 Subject Classification: 06F25
ARCIDMEDEAN SEMI-GROUP - 1) A totally
ordered semi-group all strictly-positive (strictly-
negative) elements of which belong to the same
Archimedean class. All naturally ordered Archimedean
semi-groups (cf. Natura1ly ordered groupoid) are iso-
239
ARCHIMEDEAN SEMI-GROUP
morphic to some sub-semi-group of one of the follow-
ing semi-groups: the additive semi-group of all non-
negative real numbers; the semi-group of all real
numbers in the interval (0, 1) with the usual order and
with the operation ab = min {a + b, I}; the semi-group
consisting of all real numbers in the interval (0, 1) and
the symbol 00 with the usual order and with the opera-
tions: a+b if a+b.;;;l,
ab ={
00 ifa+b>l.
The former case occurs if and only if S is a semi-group
with cancellation.
References
[1] FUCHS, L.: Partially ordered algebraic systems, Pergamon, 1963.
a.A. Ivanova
2) A semi-group S which satisfies the following con-
dition: For any a, bES there exists a natural number n
such that an ESbS. If an ESb (an EbS), the semi-group
S is called left (right) Archimedean. For commutative
semi-groups all these concepts are equivalent. Any
commutative semi-group S is uniquely decomposable
into a band of Archimedean semi-groups (and such a
decomposition coincides with the finest decomposition
of S into a band of semi-groups). This result may be
generalized in a different manner to non-commutative
semi-groups [1]. A semi-group S with an idempotent is
Archimedean (right Archiroedean) if and only if it has
a kernel K and if K contains an idempotent (K is a
right group, cf. also Kernel of a semi-group) while the
Rees quotient semi-group (cf. Semi-group) is a nil
semi-group. Archimedean semi-groups without idempo-
tents are harder to study. A complete description in
terms of certain constructions, which is especially clear
for semi-groups with a cancellation law [2], [3], was
given for the commutative case only.
References
[1] PurCHA, M.S.: 'Band of t-Archimedean semigroups', Semi
group Forum 6 (1973), 232-239.
[2] CLIFFORD, A.H. and PRESTON, G.B.: The algebraiC theory of
semigroups, 1-2, Amer. Math. Soc., 1961-1967.
[3] TAMURA, T.: 'Construction of trees and commutative
Archimedean semigroups', Math. Nachr. 36, no. 5-6 (1968),
255-287.
L.N. Shevrin
AMS 1980 Subject Classification: 06F05, 20MXX
ARCHIMEDEAN SPIRAL - A plane transcendental
curve the equation of which in polar coordinates has
the form: p = acp.
240
It is described by a point M moving at a constant rate
along a straight line d that rotates around a point a
lying on that straight line. At the starting point of the
motion, M coincides with the centre of rotation a of
the straight line (see Fig.). The length of the arc
between the points M I (PI, <1>1) and M 2(P2, <1>2) is
1= f [cp~ +In(cp+ ~)]:.
The area of the sector bounded by an arc of the
Archimedean spiral and two radius vectors PI and P2,
corresponding to angles <1>1 and <1>2, is
3_ 3
S=~
a
An Archimedean spiral is a so-called algebraic spiral
(cf. Spirals). The generalization of the Archimedean
spiral is called a neoid, the equation of which in polar
coordinates is
p = acp+l.
The spiral was studied by Archimedes (third century
B.c.) and was named after him.
References
[1] SAYELOV, A.A.: Plane curves, Moscow, 1960 (in Russian).
D.D. Sokolov
Editorial comments.
References
[A1] LOCKWOOD, E.H.: A book of curves, Cambridge Univ. Press,
1961.
AMS 1980 Subject Classification: 53A04, 14H45,
51N10
ARCSINE DISTRIBUTION - A probability measure
on the real line whose density is zero outside the inter-
val (0,1) and is Cy'x(l-x-I /'IT if O<x<1. The
corresponding distribution function is equal to
(2/ 'IT) arcsin \Ix for O';;;;;x';;;;; 1.
The generalized arcsine distribution is employed
together with the arcsine distribution. To the general-
l
ized arcsine distribution corresponds the distribution
function Fa(x) with density
Sin'lT'lTa x-u(l-x)u-I ifO<x<l,
fu(x) =
o ifx';;;O,x~l,
if 0<0:<1. The density !1/2(X) coincides with the den-
sity of the arcsine distribution. The generalized arcsine
distribution is a special case of the beta-distribution.
The first-order moment of the generalized arcsine dis-
tribution is 1-0:, and its variance is (1-0:)0:/2. The
arcsine distribution and the generalized arcsine distri-
bution occur in the study of the fluctuations of random
walks, in renewal theory (cf. Arcsine law), and are used
in mathematical statistics as special cases of the beta-
distribution.
 AREA

References lim !11-

L() - l for 0 <y<oo.
[1) FELLER, W.: An introduction to probability theory and its appli x-+oo x
cations, 1-2, Wiley, 1957-1971.
There exists a close connection between the arcsine law
[2) KENDALL, M.G. and STUART, A.: The advanced theory of statis
tics. Distribution theory, Griffin, 1969. in renewal theory and the arcsine law governing a ran-
B.A. Rogozin dom walk [3].
AMS 1980 Subject Classification: 60EXX
References
[1) FELLER, W.: An introduction to probability theory and its appli-
ARCSINE LAW - A limit theorem describing the cations, 2, Wiley, 1957-1971.
fluctuations of a random walk on the real line, which [2) SPITZER, F.: Principles of random walk, Springer, 1976.
[3) ROGOZIN, B.A.: 'The distribution of the first ladder moment
results in an arcsine distribution or a generalized and height and fluctuation of a random walk', Theory Probab.
arcsine distribution. The following feature of a Appl. 16, no. 4 (1971),575-595. (Teor. Veroyatnost. i Primenen.
Brownian motion {~I: t~O, ~o=O} was noted in 1939 16, no. 4 (1971), 593-613)
B.A. Rogozin
by P. Levy. Let TI be the Lebesgue measure of the set
{u: ~u >0, O",;:;u"';:;t} or, in other words, the time spent Editorial comments. The function L in the article above
is called a slowly varying function, cf. [1], p. 269.
by a Brownian particle on the positive semi-axis during
the interval of time [0, t]. The ratio TI / t will then have AMS 1980 Subject Classification: 60FXX
the arcsine distribution:
AREA - A numerical characteristic assigned to
p{ ~ <x} = F I/ 2(X) = ; arcsin Yx, O";;;x,,;;;l, t>O. planar figures of a certain class (for example, polygons)
and having the following properties: 1) the area is
It was subsequently noted [2] that a random walk with non-negative; 2) the area is additive (in the case of
discrete time obeys the following arcsine law: Let polygons this means that if a figure P U Q is com-
So =0, ... ,Sn, ... , be the successive locations in the posed of two figures P and Q not having common inte-
random walk, n rior points, then area(P U Q) = area P + area Q); 3) the
Sn = ~~u, area is preserved under displacements; and 4) the area
u=1
of the unit square is 1. The term 'area' is used also in a
where ~l, b ... , are independent and identically dis-
more general sense as a numerical characteristic of
tributed, let Tn be equal to the number of indices k
two-dimensional surfaces in three-dimensional space,
among 0, ... ,n for which Sk >0, and let
for k-dimensional surfaces in an n-dimensional
Kn = min {k: Sk = max Sm}, Euclidean or Riemannian space (2",;:;k"';:;n), and for the
O~m~n
boundaries of sets and other objects, see below.
then the relationships The area of a planar figure. Historically, the area was
determined initially on the class of polygons (figures
limp{~<x}
n
n ______ oc
= lim
n~oo
p{Knn <x} = Fa(x), having a decomposition into a finite number of trian-
gles without common interior points). It is important
lim P{SI <O}+ ... +P{Sn<O} =a that an area with the properties 1) - 4) exists and is
n ______ oo n
unique in the class of polygons [1], [2]. One of the
are all satisfied or not satisfied at the same time; here, consequences of 1) - 4) is that the area of an entire fig-
Fa(x) for 0<0:< 1 is the generalized arcsine distribu- ure is not less than the area of its parts.
tion, In Antiquity, the existence and uniqueness of the
FI(x) = E(x) and Fo(x) = E(x-l), area with properties 1) - 4) was assumed without an
where E(x)=O if x",;:;o and E(x)= 1 if x>O. explicit description of the class of figures; attention was
The arcsine law in renewal theory states that for concentrated on ways of calculating the area. The for-
0<0:< 1 the following equalities are valid: mula for the area of a rectangle, including one with
P{~I;:;;'O} = 1 irrational sides, was based on the method of exhaustion
(d. Exhaustion, method of). The area of a triangle or
and for
y, = t-S~" of any polygon was calculated as the area of the rectan-
where 1/1 is defined by the relation S.,." <t",;:;S"",+I' gle such that the given triangle or polygon and the rec-
tangle are both built up from the same congruent parts.
limp{Y; <x} = Fa(x) It has been shown [2] that any polygons of equal area
are decomposable into the same congruent parts.
if and only if Later, the class of squarable (Jordan-measurable) fig-
ures was distinguished. A figure M in the plane is said
for x >0, where L (x) is a function which is defined for to be squarable if for any (>0 there exist polygonal fig-
x >0 and which has the property ures P and Q such that P eM e Q and

241
AREA
(areaQ-areaP)<t:. The class of squarable figures is
extremely rich. It includes, in particular, all bounded
domains in the plane with piecewise-smooth boun-
daries. However, there are also non-squarable planar
figures. In the class of squarable figures, an area with
the properties I) - 4) exists and is unique [2].
Historically, before the class of squarable figures was
considered, people knew how to calculate the areas of
some of them: the disc, a circular sector and segments,
various kinds of choral figures, and curvilinear trapezia.
These calculations were based on the method of
exhaustion by polygons. In some cases, Cavalieri's prin-
ciple (d. Cavalieri principle) was used as the basis for
such calculations. It states that if two planar figures of
this type intersect every line parallel to a given one in
segments of identical length, then these figures are
equal in area. Techniques in integral calculus (see, for
example, [3]) give convenient methods of calculating the
area of any planar domain with a piecewise-smooth
boundary. The integral calculus justifies Cavalieri's
principle.
Attempts to extend the concept of area to more gen-
eral planar sets while retaining the properties 1) - 4) led
to measure theory and to the distinction of the class of
planar Lebesgue-measurable sets. The transition to even
more general classes of sets in the plane leads to non-
unique measures with the properties 1) - 4).
Oriented area. If there is a directed closed curve I in
the oriented plane, possibly with self-intersections and
overlaps, then for each point in the plane not lying on I
there is an integer function (positive, negative or zero),
called the degree of the point with respect on I. It indi-
cates how many times and in what direction the con-
tour I encircles that point. The integral over the entire
plane of this function, if it exists, is called the oriented
area enclosed by I. The latter differs from the ordinary
area in having a sign. See [4] for simple properties of
an oriented area.
The area of surfaces. The areas of polyhedral sur-
faces are first and most simply defined as the sums of
the areas of the planar faces. There are difficulties in
attempts to introduce the concept of area for curved
surfaces as the limit of the areas of inscribed polyhedral
surfaces (in a similar was as the length of a curve is
defined as the limit of inscribed polygonal lines). Even
for a very simple curved surface, the areas of inscribed
polyhedra with successively smaller faces may have dif-
ferent limits, depending on the choice of the sequence
of polyhedra. This is clearly demonstrated by Schwartz'
example, in which sequences of inscribed polyhedra
with different limits of the area are constructed for the
lateral surface of a right circular cylinder [2].
Most frequently, the area of surfaces is defined for
the class of piecewise-smooth surfaces with piecewise-
242
smooth boundaries (or without boundary). Usually,
this is based on the following construction. The surface
is split up into small parts with piecewise-smooth boun-
daries. In each part one selects a point at which the
tangent plane exists and projects the part orthogonally
on the tangent plane to the surface at that point; the
areas of the resulting planar projections are summed,
and finally one passes to the limit for increasingly fine
subdivisions (such that the largest of the diameters in
the parts of the subdivision tends to zero). This limit
always exists in the class of surfaces, and if the surface
is given parametrically by a piecewise C I-smooth func-
tion r( u, v), where the parameters u and v vary in a
domain D in the (u, v)-plane, then the area F is
expressed by the double integral
(1)
where gil =r~, gI2 =ru r.. g22 =r;, while ru and rv are
the partial derivatives with respect to u and v. The
proof is given, for example, in [3], [5]. In particular, if
the surface is the graph of a C I-smooth function
z = f(x,y) over a domain D in the (x,y)-plane, then
F = \l'1+.fx +.Ii dx dy. (2)
Here (1 +f; +/})I /2 = cos a., where a. is the (acute)
angle between the normal to the surface and the z-axis.
From (1) and (2) one gets the standard formulas for the
area of a sphere or parts of it, as well as methods for
calculating the areas of surfaces of rotation, etc.
For two-dimensional piecewise-smooth surfaces in
Riemannian manifolds [6], formula (1) serves as a
definition of the area, where the roles of gil, gI2 and
g22 are played by the components of the metric tensor
restricted to the surface itself.
It is important that even in the case of a two-
dimensional surface, the area is assigned not to the set
of points but to the equivalence class of mappings
representing the two-dimensional manifold into a space,
and thus differs from a measure.
k-dimensional area. For a piecewise-smooth immer-
sion f: M ~Rn of a k-dimensional manifold (with or
without boundary) into n-dimensional Euclidean space,
2:s;:;. k < n, the area is defined by means of a construc-
tion that is completely analogous to that described
above for piecewise-smooth surfaces. The only differ-
ence is that the projection now involves k-dimensional
tangent planes and that the summation is over the k-
dimensional projection volumes. If one has introduced
coordinates u I, ... , uk in a domain D eM, then the
area F of the immersion (D'!I D) is expressed by the
integral
F = ydet(g'j) du l ... du k , (3)
where gij are the scalar products (af / au i , af / au}). If

k =n -
1 and D is a domain in the coordinate hyper-
plane (x I, ... ,Xn -I), while f has the explicit form
Xn =z(xJ, ... ,Xn -d, (3) becomes
az ]2 dx
F=[ 1+ ~
n -I [
;=1
-a.
X,
l dx n -
If f is a smooth mapping, then the gij are the coeffi-
cients of the metric tensor induced by the immersion,
and it follows from (3) that the area defined by the
external construction belongs to the intrinsic geometry
of the immersed manifold. In general, (3) is taken as
the definition of area in the case of an immersion not
only into Rn but also into a Riemannian manifold.
On the class of piecewise-smooth immersed mani-
folds, the area is: a) non-negative; b) equal to 1 on the
k-dimensional unit cube in Rn; c) preserved under
orthogonal transformations; d) additive; e) semi-
continuous, i.e. F({):S;;;;liminfi-+OO F(f;) if J;~f uni-
formly. Further: f) if </>: Rn~Rn is a non-stretching
mapping, then F(</>0j):S;;;;F(j); and g) if {Pi} is a set of
(V pairwise-orthogonal k-dimensional planes and if Pi
is the projection onto Pi, then
(;)
F(p;0f> .;;;; F(j)';;;; ~F(P;0f>.
;=1
Further extensions. The theory of area. The extension
of F to more general objects while retaining some of
properties a) - g) is possible in various ways and leads
to various results. The theory of area deals with these
extensions. See [10] for a review of the relations
between various concepts of area.
The boundary between the theory of area and meas-
ure theory is not clear-cut. Traditionally, the theory of
area relates primarily to the study of the areas of con-
tinuous mappings, where multiplicity is taken into
account and the preservation of additivity is less impor-
tant. See, for example, Hausdorff measure; Favard
measure or k-dimensional measures in n-dimensional
space.
A, area can be introduced on the basis of an
approximating, integral-geometric, or functional
approach or else axiomatically. The most common con-
cepts are given below.
The Lebesgue area. It (cf. [7], [9]) is defined by
L(M,f) = lim F(j;),
where M is a finitely-triangulable k-dimensional mani-
fold, J;: M ~Rn are all possible sequences of piecewise-
linear mappings with J;~f, and F(f;) is the k-
dimensional area of the corresponding polyhedral sur-
face. The Lebesgue area is the same for Frechet-
equivalent mappings and therefore is a characteristic of
a Frechet surface.
If k = 2, the condition L(M,f) < 00 implies various
AREA
useful properties of the surface (for example, the possi-
bility of introducing isothermal parameters); in that
case, the Lebesgue area is a convenient tool sufficient
to solve the Plateau problem and more general two-
I (4) dimensional problems in variational calculus. There
arise difficulties in the investigation of variational prob-
lems for the class of continuous mappings for k>2 (the
compactness problem), which has led to the search for
other objects (currents and varifolds) and related
characteristics of area type. The use of the Lebesgue
area is restricted also by the complexity in establishing
its relationships to other concepts of area and certain
of its properties (for example, the right-hand inequality
in (5)). There are two notable features of L(M,f): for
L(M,f) =0 it may happen that the volume
V({(M)) >0, and, secondly, even for k =2, decompos-
ing Minto M 1 and M 2 with a common boundary in
the form of a curve, it may happen that
L(M,j) > L(M lo fIM,)+L(M 2 ,fIM,). (7)
Studies have also been made concerning which of the
properties a) - g) can be retained when additivity is
replaced by semi-additivity, e.g. assuming an inequal-
(5) ity of the type of (7) one still gets the Lebesgue area.
Topics of this type have not yet (1983) been examined
completely [7].
Integral-geometric areas. (Cf. [9].) Let N(/, </>,y) be
some multiplicity function for mappings </>: M~Rn at a
point YERk, for example N =card </>-I(y). Then for
f: M ~Rn, one may define the integral-geometric area
f (N(M,Paf,y)dV(y)dv(o), (8)
G(n, k)A.k
where G(n, k) is the Grassmann manifold of k-
dimensional subspaces Rk eRn, v is the normalized
Haar measure on G(n, k) and pa are the orthogonal
projections on eJEG(n, k). The areas (8) may differ for
various multiplicity functions. A very special choice of
N(M, </>,y) makes (8) coincide with the Lebesgue area
for k =2 for triangulable M, as well as for k >0 if the
Hausdorff measure Hk + 1({(M)) vanishes completely
[9]. For the case k =2, various integral-geometric areas
have been introduced also by G. Peano, Hetz and S.
Banach [7].
The area of the boundary of a set. The Minkowski
(6) area was introduced in connection with the proof of the
Brunn - Minkowski inequality and the classical isoper-
imetric inequality. It is assigned to a set A eRn, but
characterizes the area of the boundary, and is defined
by
F(a) = lim hI (V(A +hB)- V(A, (9)
h~O
where B is the unit sphere in Rn and V denotes the
volume. For sets A with a piecewise-smooth boundary
and for convex sets A, the ordinary limit in (9) exists; it
243
AREA
coincides with the (n -I)-dimensional area of the
boundary. Definition (9) still applies to sets in finite-
dimensional normed spaces, where the area (9) may
differ from the Hausdorff measure Hn - I (aA) even for
convex sets A.
There is another useful characteristic analogous to
area, which is put into correspondence with a set, but
which characterizes mainly the boundary, namely the
perimeter of a measurable set. This is a particular case
of the concept of a current mass.
Intrinsic area. If M is metrized and f: M ~Rn is a
locally isometric mapping, one has to consider the rela-
tion between L(M,f) and the Hausdorff measure
Hk(M). When M is a two-dimensional manifold of
bounded c1ln'ature, one has L(M,f)=H 2 (M). In gen-
eral, on the other hand, a continuous f: M ~Rn induces
a generalized metric p distinguished by the possibility
p(x,y)=oo on the connected components f-I(u),
u ERn. The construction of the k-dimensional Haus-
dorff measure applied to p gives a characteristic that
can be taken as the intrinsic area of an immersion. For
fin a Lipschitz class, this coincides with L(M,f) [11].
Current masses and varifolds. The integration of k-
forms with respect to a k-dimensional piecewise-smooth
oriented manifold M imbedded in Rn leads to the
current TMq,= t(q,(x), v(xdF(x), being a linear func-
tional of the k-forms q, on Rn. Here p is a unit k-vector
tangent to M. The linear functional TM essentially
characterizes M. Also, a non-linear functional can be
defined (in the same way and for non-oriented M): the
varifold VMq,= 1M I (q" p) I dF. The integral norm
(masses) II TM II and II VM I coincide with the area, i.e.
with F(M). The inclusion of the class of piecewise-
smooth submanifolds of Rn in more general current
and varifold classes plays the same role in variational
calculus as do generalized solutions in the theory of
partial differential equations.
There are extensive classes of integer currents and
varifolds among all of these. These retain many
geometric properties of submanifolds. For example, an
integer-valued current is a current having a representa-
tion T = ~;: I ni TM" where ni are integers and Mi are
C I-smooth submanifolds, subject to the condition that
the mass of the (k - I )-dimensional current dT defined
by dT(q,) = T(dq,) is finite. The masses of integer-valued
currents and varifolds can be considered as an exten-
sion of the concept of the area of a surface. Here again,
there is a relationship to the Lebesgue area. Let
piecewise-smooth mappings j;: M ~Rn converge uni-
formly, j;~f, and let L(M,f) = limi->ooL(M,j;). Then
the corresponding varifolds Vi, converge weakly to
some integer-valued varifold Vf' where
I Vf I = L(M, f). Thus, each f: M ~Rn with
244
L(M,f)<oo is naturally put into correspondence with
a varifold Vf with mass L(M,f); see [13] for this in the
language of currents.
References
[1) LEBESGUE, H.: 'Sur la mesure des grandeurs', Monogr. de
I'Enseign. Math. 1 (1935).
[2) Encyclopaedia of elementary mathematics, Vol. 5, Moscow, 1966
(in Russian).
[3) IL'IN, V.A. and POZNYAK, E.G.: Foundations of mathematical
analysis, 1, Moscow, 1971 (in Russian).
[4) LoPSHITS, A.M.: Calculating the area of oriented figures, Mos-
cow, 1956 (in Russian).
[5) POGORELOV, A.V.: Differential geometry, Noordhoff, 1959
(translated from the Russian).
[6) RASHEWSKl, P.K. [P.K. RAsHEvsKIi): Riemannsche Geometrie
und Tensorana/yse, Deutsch. Verlag Wissenschaft., 1959
(translated from the Russian).
[7) CESARI, L.: Surface area, Princeton Univ. Press, 1956.
[8) FEDERER, H.: Geometric measure theory, Springer, 1969.
[9) FEDERER, H.: 'Measure and area', Bull. Amer. Math. Soc. 58,
no. 3 (1952), 306-378.
[10) BURAGO, Yu.D. and ZALGALLER, V.A.: Geometrical inequali-
ties, Leningrad, 1980 (in Russian).
[II) BUSEMANN, H.: 'Intrinsic area', Ann. of Math. Ser. 248(1947),
234-267.
[12) ALMGREN, F.J.: The theory ofvarifolds, Princeton Univ. Press,
1965.
[13) FEDERER, H.: 'Currents and area', Trans. A mer. Math. Soc. 98,
no. 2 (1961), 204-233.
Yu.D. Burago
V.A. Zalgaller
L.D. Kudryavtsev
AMS 1980 Subject Classification: 53A05, 53825,
53C65, 51 M25, 26815, 28A25, 49F20, 58A22,
58A25
AREA FUNCflON - The set function on a sphere Q
equal to the area S (E) of that part of the convex sur-
face F that has spherical image E C Q. This definition
remains meaningful for general convex surfaces and it
gives a totally-additive set function on the ring of Borel
sets.
References
[I) ALEKSANDROV, A.D.: Mat. Sh. 3, no. I (1938),27-44.
[2] BUSEMAN, H.: Convex surfaces, Interscience, 1958.
M.I. Voftsekhovskif
Editorial comments. In the article, Q is the unit sphere in
R3 with centre at the origin. If one associates to each point
x E F the unit normal vector nx and shifts this vector nx so
that its base coincides with the origin, then the end pOint of
nx is a point x on Q The point x' is called the spherical
image of x. The procedure for obtaining the spherical image
of a point goes under the name of spherical mapping, cf.
Spherical map.
AMS 1980 Subject Classification: 53A05, 52AXX
AREA-FUNCTION - A function inverse to a hyper-
bolic function, i.e. one of the functions: arcsinh,
arccosh, arctanh, arccotanh. See also: Imerse hyper-
bolic functions.
 AREA PRINCIPLE

Editorial comments. functions f(z), z EB, with B a domain, one usually

References understands any inequality having the property that
[1] ABRAMOWITZ, M. and STEGUN, l.A.: Handbook of mathemati- equality holds if and only if the area of the complement
cal functions, National Bureau of Standards, 1964. Dover G of f(B) is zero, and the same applies for a class of
reprint, 1972.
systems of univalent functions {/k(z): z EBdZ =0,
AMS 1980 Subject Classification: 26A09 n = 1,2, ... , where Bk is a domain and G is the com-
plement of U ~ =o!k(Bj ). Usually, such a theorem is
AREA METHOD - A method of solving problems in proved by means of the area principle. That is, one
the theory of univalent functions by means of area considers any regular function Q(w), or more generally,
theorems (cf. Area principle). one having a regular derivative, on G and calculates the
Editorial comments. area a(Q(G of the image of G under the mapping of
References the function Q. Therefore, (2) is a certain extremely
[A 1] DUREN, P.L.: Univalent functions, Springer, 1983, Chapt. 4. general area theorem in the class ~.
Let FE~ and let
AMS 1980 Subject Classification: 30C50
1n F(t)- F(z) = ~
"'" Wp,q
t -p
z
-q
, t, z E~ .
A'

t -z p,q~1
AREA PRINOPLE - The area of the complement to
the image of a domain under a mapping by a function One chooses a suitable function, regular in CF(I::.,), to
regular in it is non-negative. The area principle was be able to write (2) as
first used in 1914 by T.H. Gronwall [1], who in this
way proved the so-called exterior area theorem for func- q~lq Ip~lwp,qXp 12 .;;; P~I~ I Xp 1
2
,
(3)

tions of class ~ - functions

where the xp are any numbers not simultaneously equal
al
F(z) = z +Ilo +-+
z
... , to zero and such that
lim I xp II Ip < 1.
that are regular and univalent in the annulus /:l: = {z: p ..... oo

1< I z 1< oo} (cf. Univalent function). The area More general area theorems in ~ have also been
a(CF(/:l' of the complement of the image F(I::.') of t::.' obtained [5).
under a mapping W=F(Z)E~ can be determined by the Area theorems have been proved for the following:
formula the class IDC(alo ... ,an) of systems {/k(z): /k(O)=ako
z EI::. H= I of functions /k that map the disc I::. confor-
(J(CF(~' = 'TT [1- k~lk I ak 12] ;;;. mally and univalently onto domains without pairwise
common points, i.e. onto non-adjacent domains [6]; the
and consequently class ~(B) (~(B), B 300, is the class of functions F that
00
~klakI2.;;;1. (1) are regular and univalent in B \ { oo} and such that
k ~I F( 00) = 00, limz--+oo F(z) / z = 1), [7]; and non-
By means of (1) the first results were obtained for func- overlapping multiple-connected domains (see [6] and
tions of the classes ~ and S, where S is the class of also [8], [9]). All the area theorems for multiple-
functions f(z)=z + ~:=2akzk that are regular and connected domains can be proved by the method of
uni" lent in the disc I::.={z: I z I <I} (cf. Bieberbach contour integration (cf. Contour integration, method ot).
conjecture; Distortion theorems). A more general area By the area method one understands methods of solv-
theorem has been proved [2]. G.M. Goluzin [3] ing various problems in the theory of univalent func-
extended the area theorem to p-valent functions in the tions by the use of area theorems.
disc (cf. Multivalent function). For instance, from (3) one may obtain by means of
The following area theorem has been proved [4]: Let the Cauchy inequality that
!:. E~, E=CF(I::.'), and let Q(w) be a regular function in
B.1f I ~ wP.qXpX~ 12 .;;; ~ 1- I
~q~1 p~IP
Xp 12 ~ 1- I X~ 12,
q~lq
(4)
+00
Q(F(z)) = ~ ak z -k 'I< const; where xp and x~ are such that the series on the right-
k~-oo
hand side converge. If in (4), for example, xp = t -P,
then +00 x~ =z -q, I t I > 1, I z I> 1, one obtains the chord-
~ k I ak 12 .;;; 0, (2) distortion theorem
k~-oo

and equality holds only if the area a(E) of E is zero. In F(t) - F(z) 12 .;;; In I t 12 In I z 12 .
l t-z It1 2-1 Iz12-1
By the area theorem for a certain class of univalent

245
AREA PRINCIPLE
Area theorems, for example in the class WC(a I, . . . , an),
give necessary and sufficient conditions for a system
{}k(z): jk(O)=ak, ZELl}Z=1 of meromorphic functions
jk to belong to the class WC(al,' .. ,an) (see [6], p.
179).
References
[1] GRONWALL, T.H.: 'Some remarks on conformal representa-
tion', Ann. of Math. Ser. 2 16 (1914 - 1915),72-76.
[2] PRAWITZ, H.: 'Ueber Mittelwerte analytischer Funktionen',
Arkiv. Mat. Astron., Fysik lOA, no. 6 (1927), 1-12.
[3] GOLUZIN, G.M.: 'On p-valued functions', Math. Sb. 8, no. 2
(1940), 277-284 (in Russian).
[4] LEBEDEV, N.A. and MILIN, I.M.: 'On the coefficients of certain
classes of analytic functions', Mat. Sb. 28, no. 2 (1951),359-
400. (in Russian).
[5] NEHARI, Z.: 'Inequalities for the coefficients of univalent func-
tions', Arch. Rational Mech. and Anal. 34, no. 4 (1969),301-
330.
[6] LEBEDEV, N.A.: The area principle in the theory of univalent
functions, Moscow, 1975 (in Russian).
[7] MILIN, I.M.: Univalent functions and orthonormal systems, Mos-
cow, 1971 (in Russian).
[8] ALENITSYN, YU.E.: 'Area theorems for functions that are ana-
lytic in a finitely-connected domain', Izv. Akad. Nauk SSSR
Ser. Mat. 37, no. 5 (1973), 1132-1154 (in Russian).
[9] GULTYANSKIi, V.YA. and SHCHEPETOV, V.A.: 'A general area
theorem for a certain class of q-quasi-conformal mappings',
Dokl. Akad. Nauk SSSR 218, no. 3 (1974), 509-512 (in Rus-
sian).
[10] GRUNSKY, H.: 'Koeffizientenbedingungen fUr schlichtabbil-
dender meromorphe Funktionen', Math. Z. 45, no. 1 (1939),
29-61.
N.A. Lebedev
Editorial comments. The inequalities (3) and (4) are a
form of the Grunsky inequalities.
References
[A 1] DUREN, P.L.: Univalent functions, Springer, 1983, Chap!. 1 O.
AMS 1980 Subject Classification: 30C50, 30CXX
ARF-INV ARIANT, invariant of Arf - An invariant of
a quadratic form modulo 2, given on an integral lattice
endowed with a bilinear skew-symmetric form. Let II
be an integral lattice of dimension k = 2m and let
~(x,y) be a form for which ~(x,y)= -~(y, x). There
exists bases {e I ,II, ... , em,fm }, called symplectic
bases, in which the matrix of ~(x, y) reduces to block-
diagonal form: The diagonal contains the blocks
l.e.
IHei ,!,) = -tJ;(j;, e,l = 1,
while the other entries are zero.
Suppose that a mapping
tJ;0: IT ..... Z2
is given on II such that
tJ;o(x+y) = tJ;o(x)+tJ;o(Y)+tJ;(x,yl Il)od2
(a 'quadratic form modulo 2'). The expression
m important geometric principles of function theory, e.g.
~ tJ;o(e,)tJ;o(j;)
I I
246
is then called an Arf-invariant [1]. If this expression
equals zero, then there is a symplectic basis on all ele-
ments of which ~o vanishes; if this expression equals
one, then there is a symplectic basis on all elements of
which, except e I and fl, the form ~o vanishes, while
tJ;o(ed = tJ;O(jI) = l.
References
[1] ARF, c.: 'Untersuchungen tiber quadratischen Formen in
K6rpem der Charakteristik 2, 1', J. Reine Angew. Math. 183
(1941), 148-167.
A. V Chernavskif
Editorial comments. See [A 1], Appendix 1, for material
concering the Arf-invariant for inner product spaces over a
field F of characteristic 2 in relation to the Witt algebra of
quadratic inner product spaces over F.
References
[A1] MILNOR, J. and HUSEMOLLER, D.: SymmetriC bilinear forms,
Springer, 1973.
AMS 1980 Subject Classification: 10C05
ARGUMENT - 1) The argument of a function is the
variable (also called the independent variable) on which
the value of the function depends.
2) The argument of a complex number
Z = x + ry = r( cos <p + i sin <p), represented in the plane by
the point with coordinates x and y, is the angle <p of the
radius vector r of this point with the x-axis.
AMS 1980 Subject Classification: 30-01
ARGUMENT, PRINOPLE OF THE, argument prin-
ciple - A geometric principle in the theory of functions
of a complex variable. It is formulated as follows: Let
D be a bounded domain in the complex plane C, and
let, moreover, the boundary aD be a continuous curve,
oriented so that D lies on the left. If a function w = f(z)
is meromorphic in a neighbourhood of D and has no
zeros or poles on aD, then the difference between the
number of its zeros N and the number of its poles P
inside D (counted according to their multiplicity) is
equal to the increase of the argument of f when travel-
ling once around aD, divided by 27T, i.e.
I
N- P = 27T !:J.aDargj,
where argf denotes any continuous branch of Argf on
the curve aD. The expression on the right-hand side
equals the index indof(aD) of the curve f(aD) with
respect to the point w =0.
The principle of the argument is used in the proofs
of various statements on the zeros of holomorphic func-
tions (such as the fundamental theorem of algebra on
polynomials, the theorem of Hurwitz on zeros, etc.).
From the principle of the argument follow many other
the principle of invariance of domain (cf. Invariance of

domain, principle of), the maximum-modulus principle
and the theorem on the local inverse of a holomorphic
function. In many questions the principle of the argu-
ment is used implicitly, in the form of its corollary: the
Rouche theorem.
There are generalizations of the principle of the argu-
ment. The condition that f be meromorphic in a neigh-
bourhood of 15 may be replaced by the following: f has
only a finite number of poles and zeros in D and
extends continuously to aD. Instead of the complex
plane, an arbitrary Riemann surface may be con-
sidered: the boundedness of D is then replaced by the
condition that 15 be compact. From the principle of the
argument for a compact Riemann surface it follows
that the number of zeros of an arbitrary meromorphic
function, not identically equal to zero, is equal to the
number of poles. The principle of the argument for
domains in C is equivalent to the theorem on the sum
of the logarithmic residues (d. Logarithmic residue).
For this reason, the following statement is sometimes
called the generalized principle of the argument. If f is
meromorphic in a neighbourhood of a domain 15 which
is bounded by a finite number of continuous curves
and if f has no zeros or poles on aD, then for any fun::
tion cp which is holomorphic in a neighbourhood of D
the equality
1 1" {z I N P
-2. (<j>(Z)T(
) dz = ~<j>(ad- ~<j>(bk)
'lTl iD J' Z k= I k- I
holds, where the first sum extends over all zeros and
the second sum extends over all poles of finD. There
is also a topological generalized principle of the argu-
ment: The principle of the argument is valid for any
open mapping f: D~C that is locally finite-to-one and
extends continuously to aD, while 0, 00 fI-f(aD).
An analogue of the principle of the argument for func-
tions of several complex variables is, for example, the
following theorem: Let D be a bounded domain in
with Jordan boundary aD and let f: 15~c,,- be a holo-
morphic mapping of a neighbourhood of D such that
Otlf(aD); then the number of pre-images of 0 in D
(counted according to multiplicity) is equal to
indof(aD).
References
[I] LAWRENTJEW, M.A. [M.A. LAVRENT'EV] and SHABAT, B.V.:
Methoden der komplexen Funktionentheorie, Deutsch. Verlag
Wissenschaft., 1967 (translated from the Russian).
[2] SHABAT, B.V.: Introduction to complex analySiS, 2, Moscow,
1976 (in Russian).
E.M. Chirka
AMS 1980 Subject Classification: 30AXX, 32AXX
AluTHMEnc - The science of numbers and opera-
tions on sets of numbers. Arithmetic is understood to
include problems on the origin and development of the
ARITHMETIC
concept of a number, methods and means of calcula-
tion, the study of operations on numbers of different
kinds, as well as analysis of the axiomatic structure of
number sets and the properties of numbers. When
referring to the logical analysis of the concept of a
number, the term theoretical arithmetic is sometimes
employed. Arithmetic is closely connected with algebra,
which includes the study of operations performed on
numbers. The properties of integers form the subject of
number theory (d. Elementary number theory; Number
theory).
The term 'arithmetic' is also sometimes employed to
denote operations performed on objects of very dif-
ferent kinds: 'matrix arithmetic', 'arithmetic of qua-
dratic forms', etc.
The art of computation arose and developed long
before the times of the oldest written records extant.
The oldest mathematical records are the Cahoon papyri
and the famous Rind papyrus, which is believed to date
back to about 2000 B.c.. An additive hieroglyphic sys-
tem for the representation of numbers (d. Numbers,
representations of) enabled the ancient Egyptians to
perform addition and subtraction operations on natural
numbers in a relatively simple manner. Multiplication
was carried out by doubling, i.e. the factors were
decomposed into sums of powers of two, the individual
summands were multiplied, and the components added.
Operations on fractions (cf. Fraction) were reduced in
Ancient Egypt to operations on aliquot fractions, i.e. on
fractions of the type I / n. More complicated fractions
were decomposed with the aid of tables into a sum of
aliquot fractions. Division was carried out by subtract-
ing from the number to be divided the numbers
obtained by successive doubling of the divisor. The
clumsy hexadecimal system of the ancient Babylonians
gave rise to many difficulties in executing arithmetic
operations. There are numerous tables employed by
en ancient Babylonians to effect multiplication and divi-
SIOn.
In Ancient Greece arithmetic was conceived as the
study of the properties of numbers; practical calcula-
tions were not included. Problems on the technique of
operations on numbers, i.e. methods of calculation,
were considered to be an independent science - logis-
tics. This differentiation was inherited from the Greeks
by Europeans in the Middle Ages. It was only with the
advent of the Renaissance that both the beginnings of
the theory of numbers and the practice of calculation
began to be included in the concept of arithmetic.
Greek mathematics made a sharp distinction between
the concepts of a number and of a magnitude. To
Greek mathematicians numbers meant only what is
known as natural numbers in our own days (d. Natural
number); they also distinguished between concepts of
247
ARITHMETIC
different types, such as a number and a geometrical
magnitude. There are no logistic studies by ancient
Greeks extant, but it is known that their multiplication
technique was close to our own. Their alphabetic
number system severely complicated operations carried
out on numbers. In Ancient Greece ordinary fractions
were employed in calculations; however, a fraction was
not considered as a number in its own right, but merely
as a quotient of natural numbers.
Books 7 - 9 of Euclid's Elements (3-rd century B.C.)
deal exclusively with arithmetic in the sense in which
the word was employed in ancient times. They mainly
deal with the theory of numbers: the algorithm for
finding the greatest common divisor (cf. Euclidean algo-
rithm), and with theorems about prime numbers (cf.
Prime number). Euclid proved that multiplication is
commutative, and that it is distributive with respect to
the operation of addition. He also studied the theory
of proportion, i.e. the theory of fractions. Other books
of his treatise include the general theory of relations
between magnitudes, which may be considered as the
beginning of the theory of real numbers (cf. Real
number).
The manuscripts of Diophantus (probably 3-rd cen-
tury A.D.) that are left to us contain operations on
powers (not greater than 6) of numbers, and a few
examples of subtraction operations. These are impli-
citly operations on negative numbers. Diophantus
applied his own rules to rational numbers only (cf.
Rational number).
Chinese mathematicians in the 2-nd century per-
formed operations on fractions and negative numbers.
At a somewhat later period they studied the extraction
of square roots and cubic roots, the approximate values
of which were expressed as decimal fractions (cf.
Decimal fraction). The methods employed by Chinese
mathematicians in solving problems in arithmetic, in
particular the rule of two false assumptions, appears in
several manuals on arithmetic, first in Arab, then in
European manuals. Very little is known about the ini-
tial development of arithmetic in India. The simplest
fractions were utilized in India long before the Chris-
tian era. Our own decimal computation system is of
Indian origin. The earliest written Indian mathematical
records extant were compiled in the 5-th century and
indicate that the knowledge of arithmetic in India of
that period was of a high standard. Indian mathemati-
cians operated on integers and fractions using methods
very similar to our own. They solved many problems
on proportions and on the rule-of-three, and could
compute percentages. Studies on negative numbers
began in India in the 7-th century. The works of
Bhaskara II The wreath of science (l2-th century) con-
tain rules for the multiplication and division of negative
numbers.
248
Indian mathematics exerted a decisive influence on
the development of the knowledge of the arithmetic of
the Arabs. The treaty on arithmetic written by
Muhammed AI-Khwarizmi in the 9-th century greatly
contributed to the dissemination of the Indian decimal
notation and of the methods of addition, subtraction,
multiplication, division and extraction of square roots.
Many ancient nations performed their calculations
on an abacus, which had replaced the primitive count-
ing on fingers. The external form of the abacus varied,
but its principle remained the same: ruling a column,
or some other cascade-type marking of numbers. An
abacus was used by the Greeks long before the Chris-
tian era. The Chinese suan-pan and the Russian schety,
which externally resemble each other, are also variants
of the abacus.
While the European studies in the domain of number
theory are based on Greek mathematics, in particular
on the work of Euclid and Diophantus, this is not true
of computational techniques. The development of arith-
metic in Europe is connected with the arrival of the
Indian positional decimal system and of Arabic
numerals in Europe. The technique of arithmetic opera-
tions was not taken over from India directly, but rather
by way of the work of Al-Khwarizmi and other Arab
mathematicians.
The abacus was extensively employed in the Middle
Ages. It also became a synonym of the word arithmetic,
so that Leonardo da Pisa (13-th century) called his
arithmetical treatise The book of abacus. The book
quotes calculation methods taken from Arab sources,
while introducing many original improvements. Thus,
the addition of fractions is performed by way of the
least common multiple of the denominators, and the
operations are checked not only, like the Indians, in the
base nine, but also in certain other bases. The problems
dealt with in the book include the rule-of-three, the rule
of alligation, problems involving recurrent sequences,
arithmetic progressions and geometric progressions (cf.
Arithmetic progression; Geometric progression). In
Europe decimal fractions began to be employed in the
15-th century, but became widely known only in the
16-th century following the publication of the studies of
S. Stevin.
Various methods for the multiplication and division
of multiplace numbers were proposed in the 15-th, 16-
th and in the later centuries. They differed from each
other only by the system of notation of the intermedi-
ate computations. A. Riese, through numerous text-
books, was most influential in the move to replace the
old computation (in terms of counters and Roman
numerals) by the newer methods (using pen and
Hindu-Arab numerals).
In Europe negative numbers were first used by
Leonardo da Pisa, who treated them as one would treat

debts. A system of operations on negative numbers was
presented in the l6-th century by M. Stiefel. He
referred to these numbers as 'fictitious'. Proofs for the
rules of operation on negative numbers were still stu-
died as late as the l8-th century, and it is only owing
to the critical reasoning in the second half of the 19-th
century that such studies have ceased to be taken seri-
ously.
In Europe, up to the l5-th or l6-th centuries, arith-
metic operations on irrational numbers were limited to
the extraction of square roots. Nevertheless, Leonardo
da Pisa considered the problem of the approximate cal-
culation of cubic roots as well as square roots. S. dal
Ferro (around 1500) and N. Tartaglia (l6-th century)
used cubic roots in solving equations of the third
degree. There was no general treatment of operations
performed on real numbers. The concept of a real
number was assimilated in mathematics only gradually,
in connection with the development of analytic
geometry and mathematical analysis. Up to the l8-th
century, proofs of operations on irrational numbers
were limited to magnitudes expressible by radicals.
Complex numbers were encountered at various times,
beginning with Indian mathematics, in solving qua-
dratic equations. However, imaginary solutions were
discarded as non-existent. The arithmetic of complex
numbers (cf. Complex number) begins with R. Bombelli
(l6-th century), who gave formal rules for performing
arithmetic operations on such numbers. However, even
as late as the l7-th century such operations were per-
formed in a manner similar to those on real numbers,
which frequently led to errors. It is only in the l8-th
century that a precise definition of the arithmetic of
complex numbers could be given, owing to the formu-
las of A. de Moivre and L. Euler.
The idea of logarithms dates back to Archimedes (3-
rd century B.C), who compared the terms of arithmetic
and geometric progressions. Stiefel (l6-th century)
extended these progressions to the left by adding nega-
tive powers. He showed that there is a connection
between operations carried out on these series, thus
introducing the fundamental idea of a logarithm. Tak-
ing logarithms as a routine stage in calculations was
introduced by 1. Napier and 1. Burgi in the first half of
the l7-th century.
The first calculating machines were built in the l7-th
century by W. Schickard and B. Pascal, who worked
independently of each other; they were the prototypes
of modern calculators. However, it was only in the 19-
th century that such machines began to be extensively
employed in practical work. The rise of fast electronic
computers in middle of the 20-th century enhanced the
importance of performing an algorithmic operation in a
minimum number of elementary operations.
It has been considered ever since the times of Euclid
ARITHMETIC
that in order to built a theory, it is sufficient to decom-
pose it into a small number of clear, simple, primary
assumptions and to ensure that all basic postulates of
the theory can be deduced from them in a purely logi-
cal manner. It was tacitly assumed that the connection
of these fundamental principles with the real world
should be immediately perceptible.
The method of models, used to give foundations of
mathematical theories, was discovered in the 19-th cen-
tury. The use of such a method was inevitable, since
certain objects and theories studied in mathematics had
not found a real-life interpretation. These include,
mainly, complex numbers, ideals, and non-Euclidean
and n-dimensional geometries. The method of models
permitted one to reduce the consistency of one
mathematical theory to the consistency of another.
Thus, on the assumption that Euclidean geometry is
consistent, is was proved that Lobachevskii's geometry
was consistent, while the consistency of Euclidean
geometry was reduced to the consistency of the arith-
metic of real numbers.
At the end of the 19-th century, the foundation of
arithmetic seemed complete. R. Dedekind and G.
Peano, independently of each other, gave an axiom sys-
tem for the natural numbers from which all known pro-
positions of this science could be deduced. K. Weier-
strass proposed to take pairs of natural numbers as a
model for the integers and positive rational numbers.
The geometrical representation of complex numbers,
discovered by J. Argand, C Wessel and CF. Gauss, is
in essence a model for the theory of complex numbers
in the framework of the theory of real numbers.
Finally, the set-theoretic approach allowed Dedekind,
G. Cantor and Weierstrass to built the theory of real
numbers.
However, following the discovery of paradoxes in set
theory, there arose the question of how to give a foun-
dation for the arithmetic of natural and real numbers.
How can one be sure there are no paradoxes in those
branches of mathematics as well? Direct perception
does not tell one that the universe is infinite or that
matter is infinitely divisible. For this reason the ideas
of the infinite set of natural numbers and the continu-
ous number axis may be considered as not directly con-
nected with the physical world. On the other hand,
there is no simpler model of the arithmetic of natural
numbers than the theory of natural numbers itself,
while models of the theory of real numbers are based to
a considerable extent on the apparatus of set theory,
the reliability of which seems to be open to doubt.
What should be the methods and means of reason-
ing, not involving models, which would directly indi-
cate the absence of inconsistencies in a given theory -
or, in other words, which would show that a logical
reasoning based on the axioms of the theory will never
249
ARITHMETIC
yield results which contradict each other?
In the view of D. Hilbert, the reason for the appear-
ance of paradoxes in set theory is that modes of rea-
soning that are undoubtedly applicable to finite sys-
tems of objects, are illegitimately applied to infinite
families. This may be avoided if the symbols employed
are considered as objects of some new theory, and if
the logical conclusions are expressed with the aid of a
formal process. In such a case a statement in a theory
is expressed as a formula constructed out of a finite set
of symbols, while a proof is represented by a finite
chain of formulas, obtained according to certain rules
from formulas known as axioms (cf. Axiom). In this
way, according to Hilbert, it will be possible to operate
on finite rather than on infinite objects and to obtain a
reliable procedure for establishing the consistency of
any theory. Hilbert hoped that this approach would
yield, in the first place, a positive solution to the prob-
lem of the consistency of the arithmetic of natural
numbers and that it would show that the addition of
any unprovable formula in number theory to the for-
mulas of arithmetic converts this system of axioms into
an inconsistent system.
However, these hopes were smashed. K. G6del in
1931 demonstrated that formal arithmetic (cf. Arith-
metic, formal) is incomplete. He found, moreover, that
for any consistent formal system containing the axioms
of arithmetic, it is possible to find an explicit descrip-
tion of some closed formula u such that neither the for-
mula u itself nor its negation can be deduced in the
formal system.
Using this result it is possible to prove that non-
isomorphic models of formal arithmetic exist. At the
same time, the system of Peano axioms is categorical.
How is this to be explained? The system of Peano
axioms contains the induction axiom: Any natural
number has a certain property P if the number 1 has
this property and if for any natural number n that has
the property P the natural number n + 1 also has this
property. In this axiom P may stand for any conceiv-
able property of natural numbers. In the corresponding
axiom of formal arithmetic, P may denote only such
properties of natural numbers as can be expressed by
the methods of the given formalism. The difference
between these two axioms is immaterial if one discusses
theorems of elementary number theory, but is very
material in the explanation of the properties of the for-
mal theory.
G6del also showed that a consistent formal system
which includes formal arithmetic contains a formula
expressing its consistency, and that this formula cannot
be proved in this system. Thus, the consistency of such
a formal system can only be substantiated by means
that are stronger than those formalized in the system
250
itself. G. Gentzen in 1936 gave a proof of the con-
sistency of formal arithmetic, using transfinite induc-
tion up to the transfinite number Eo. One must natur-
ally inquire into the consistency of the means used in
the proof. Other approaches to the problem of con-
sistency of the arithmetic of natural numbers were also
investigated in this context.
Attempts to overcome the difficulties involved in the
foundation of the theory of real numbers played a cer-
tain role in the development of the constructive
approach in mathematics.
References
[I] The history of mathematics, Moscow, 1970-1973 (in Russian).
[2] WAERDEN, B.L. VAN DER: Ontwakende wetenschap, Noordhoff,
1957.
[3] KLEENE, S.c.: Mathematical logic, Wiley, 1967.
[4] HILBERT, D.: Grundlagen der Geometrie, Springer, 1913.
[5] Encyclopaedia of elementary mathematics, Vol. 1, Moscow-
Leningrad, 1951-1963 (in Russian).
[6] MOLODSHIi, V.N.: Fundamentals of the study of numbers in the
18-th and the beginning of the 19-th century, Moscow, 1963 (in
Russian).
[7] NECHAEV, V.l.: Number systems, Moscow, 1975 (in Russian).
[8] STRUIK, DJ.: A concise history of mathematics, Dover, reprint,
1967.
[9] THuRSTON, H.A.: The number-system, Biackie, 1956.
A.A. Bukhshtab
v.I. Pechaev
Editorial comments.
References
[A1) WAERDEN, B.L. VAN DER: Geometry and algebra in ancient
civilisations, Springer, 1983.
AMS 1980 Subject Classification: 00A25, 10-XX,
01A05
ARITHMETIC CONTINUUM - The number axis.
AMS 1980 Subject Classification: 26A03
ARITHMETIC DISTRIBUTION - A discrete proba-
bility distribution concentrated on a set of points of the
form -+-nh, where h >0 and n = 1, 2, .... It is a special
case of a lattice distribution.
A. V. Prokhorov
AMS 1980 Subject Classification: 60-XX
ARITHMETIC, FORMAL, arithmetical calculus - A
logico-mathematical calculus which formalizes elemen-
tary number theory. The language of the most common
kind of formal arithmetic contains the constant 0,
numerical variables, the equality sign, the function sym-
bols +, " ' (successor) and logical symbols. The terms
are built from the constant and the variables of the
function symbols; in particular, natural numbers are
denoted by terms of the form 0' ... '. The atomic formu-
las are equalities of terms; the remaining formulas are
built from atomic formulas by means of the logical
connectives &, V, ~, -', "d, 3. Formulas in the

language of formal arithmetic are known as arithmetical
formulas. The postulates of formal arithmetic are the
postulates of predicate calculus (classical or intuitionis-
tic, depending on the formal arithmetic in question),
the Peano axioms
a'=b' ~a=b, .....,(a'=O), (a=b&a=c)~b=c,
a=b ~ a'=b', a+O = a, a+b' = (a+b)',
a'O = 0, ab' = (ab)+a,
and the axiom scheme of induction:
A(O)& \fx(A(x)~A(x')) ~ A(x),
where A is an arbitrary formula, known as the induction
formula.
The means of formal arithmetic are adequate for the
derivation of the theorems presented in standard
courses on elementary number theory. At the time of
writing (the 1980's) it would appear that no significant
number-theoretic theorem which can be demonstrated
without recourse to analysis is known which could not
be demonstrated by formal arithmetic. To write down
and prove such theorems in formal arithmetic, the
expressive potential of the latter must be investigated.
The expression of many number-theoretic functions in
terms of formal arithmetic is especially important. In
particular, statements about primitive (and even partial)
recursive functions (cf. Recursive function) may be
rewritten in the language of formal arithmetic. In this
way it is possible to prove formulas expressing impor-
tant properties of partial recursive functions, in particu-
lar their defining equations. Thus, the equation
f(x')=g(x,f(x)) is expressed by the formula
\fa \fb \fc(F(x', a) & F(x, b) & G(x, b, c) ~ a =c),
where F and G are formulas representing the graphs of
the functions f and g, respectively. It is possible to for-
mulate statements about finite sets in the language of
formal arithmetic. Moreover, classical formal arithmetic
is equivalent to the Zermelo- Fraenkel axiomatic set
theory without the axiom of infinity: A model of any
one of these systems can also be constructed in the
other system.
The deductive force of the system of formal arith-
metic is characterized by the ordinal 0 (the least solu-
tion of the equation w =(): It is possible to deduce in
formal arithmetic the scheme of transfinite induction up
to any ordinal a<o, but not the scheme of induction
up to 0. The class of provably-recursive functions of the
system of formal arithmetic (i.e. the class of partial
recursive functions, the general recursiveness of which
can be established by the means of formal arithmetic)
coincides with the class of ordinal recursive functions
with ordinals <0. This makes it possible to extend
formal arithmetic in certain ways, e.g. to have a formal
arithmetic with symbols for all primitive recursive func-
ARITHMETIC. FORMAL
tions and the corresponding additional postulates. For-
mal arithmetic satisfies both Godel incompleteness
theorems (cf. GOdel incompleteness theorem). In partic-
ular, there exist polynomials P, Q in several variables
such that the formula 'v'X(P(X)==Q(X)) cannot be
proved even though expressing a true statement - the
consistency of the system of formal arithmetic. A com-
binatorial (Ramsey-type) theorem independent of for-
mal arithmetic has been constructed [6].
In studying the structure of a system (in particular,
the problems of consistency), formal arithmetic is for-
mulated without quantifiers, but with Hilbert's (-
symbol. The formulas of this system are quantifier-free,
but it is permitted to use terms of the type (xA (x)
(yielding some x which satisfies condition A if it exists;
otherwise 0). The postulates of the (-system are the pos-
tulates of propositional calculus, the equality rules, the
rule of substitution for the free numerical variables, the
axioms for the functions +, " ',pd (predecessor), and
the following (-axioms:
A(a) ~ A (xA (x)); A(a) ~ xA(x)=rt=a';
-A (xA (x)) ~ xA(x)=O.
The quantifiers are introduced by means of abbrevia-
tions:
3xA(x)~A(xA(x)); \fxA ~.....,3x-A.
The induction axiom turns out to be derivable.
The formulas of formal arithmetic with free variables
define number-theoretic predicates. Formulas not con-
taining free variables (closed formulas) express proposi-
tions. A k-place predicate g; about natural numbers is
called arithmetic if there is an arithmetical formula
P (x], ... ,Xk) such that for any natural numbers
n 1, . . . ,nk, the relation
P(nj, ... ,nd ~ <nj, ... ,nk>E9
is valid. This determines a classification of arithmetic
predicates by the type of prefix in the prenex form of
the respective formula. The class ~n (the class TIn) con-
sists of predicates Pea) representable by a formula of
the form
QjXj ... Qnxn(f(xj, ... ,xn, a)=O),
where f is a primitive recursive function, Q1 is 3 (or,
respectively, 'v') and Qi are alternating quantifiers (i.e.
Qi + 1 is Qi where V is 3, and "3 is 'v'). Each such class
contains a universal predicate, i.e. a predicate T(e, a)
such that for each one-place predicate P of this class
there exists a number ep for which the identity
'v'x(T(ep ' x) P(x)) is valid. Multi-place predicates can
be reduced to one-place predicates with the aid of func-
tions which code systems of natural numbers. For any
n >0 the inequality ~n==TIn is valid, and a class with a
lower index is properly included in any class with a
larger index. This classification of predicates generates
251
ARITHMETIC, FORMAL
a classification of arithmetic functions based on the
classification of their graphs. Not all number-theoretic
predicates are arithmetic: an example of a non-
arithmetic predicate is the predicate T(x) such that for
any closed arithmetical formula A the identity
T<rAl):=A is valid, where rAl is the number of the
formula A in some fixed numbering satisfying some
natural conditions. The predicate T plays an important
role in the studies of the structure of formal arithmetic,
in particular in the problem of the independence of its
axioms. The addition to formal arithmetic of the sym-
bol T with axioms of the type
T<rA&Bl) T<rAl)&TnBl), which express its com-
mutativity with logical connectives, makes it possible to
prove that formal arithmetic is consistent. The same
construction, this time inside formal arithmetic, may be
carried out for a subsystem A [n1 of formal arithmetic
in which the induction scheme is restricted by the fol-
lowing condition: The complexity of the induction for-
mula is ';;;;;n, i.e. it belongs to Iln U ~n' The function of
T is successfully performed by, say, the universal predi-
cate for ~n + I, and the corresponding proof is per-
formed in the system A [n + 11' It follows from Godel's
second incompleteness theorem that A [n + 11 is stronger
than A [n1' so that formal arithmetic does not coincide
with anyone of the systems A [n1' and the induction
scheme cannot be replaced by any finite set of axioms.
Formal arithmetic turns out to be complete with
respect to formulas from the class ~I: A closed formula
of this class is provable in formal arithmetic if and only
if it is true. Since the class LI contains an
algorithmically-unsolvable predicate, it follows that the
problem of provability in formal arithmetic is algo-
rithmically unsolvable.
If formal arithmetic is defined as the Gentzen formal
system of the usual type, a cut (the analogue of modus
ponens in Gentzen systems) cannot be eliminated.
Elimination of a cut becomes possible if the induction
scheme is replaced by the rule of infinite induction (w-
rule):
A (0) ... A(N) ...
V'xA(x)
One of the first proofs of the consistency of formal
arithmetic was obtained in this way. A manageable for-
mulation of the w-rule is: 'The number e is a number of
the general recursive function which describes the
derivation of length <io by the w-rule' (w-prooj). This
predicate belongs to the class Il 2 , and the system
obtained is not formal, but only semi-formal. To each
derivation of a formula A in formal arithmetic it is pos-
sible to assign a number e such that the statement 'e is
an w-proof of the formula A without a cut' is true, and
can even be proved in formal arithmetic. Since a cut-
free w-proof of a closed quantifier-free formula contains
252
no quantifiers, it follows that formal arithmetic is con-
sistent. If Gooel's second incompleteness theorem is
now applied, it follows that the cut elimination theorem
from any proof in formal arithmetic cannot be
demonstrated by the means of formal arithmetic;
nevertheless, for any given N a proof in formal arith-
metic is possible for any derivation with complexity of
induction ';;;;;N. Passing to w-proofs makes it possible to
establish many meta-mathematic theorems about the
system of formal arithmetic, in particular completeness
with respect to formulas of ~I and the ordinal charac-
ter of provably-recursive functions.
References
[1] KLEENE, S.C.: Introduction to metamathematics, North-Holland,
1951.
[2] HILBERT, D. and BERNAYS, P.: Grundlagen der Mathematik, 1-
2, Springer, 1968-1970.
[3] NOVIKOV, P.S.: Elements of mathematical logic, Edinburgh,
1964 (translated from the Russian).
[4] KREISEL, G.: 'A survey of proof theory', J. Symbolic Logic 33,
no. 3 (1968), 321-388.
[5] GODEL, K.: 'Ueber eine bisher nicht beniitzte Erweiterung des
finiten Standpunktes', Dialectica 12 (1958), 280-287.
[6] PARIS, J. and HARRINGTON, L.: 'A mathematical incomplete-
ness in Peano arithmetic', in 1. Barwise (ed.): Handbook of
mathematical logic, North-Holland, 1977, pp. 1133-1142.
G.E. Mints
AMS 1980 Subject Classification: 03830
ARITHMETIC FUNcrION, number-theoretic function
- A complex-valued function, the domain of definition
of which is one of the following sets: the set of natural
numbers, the set of rational integers, the set of integral
ideals of a given algebraic number field, a lattice in a
multi-dimensional coordinate space, etc. These are
arithmetic functions in the wide sense. However, the
term is often employed to denote a function of the
above type with special arithmetic properties. The most
commonly occurring arithmetic functions have tradi-
tional symbolic notations: cp(n) is the Euler function;
d (n) or T( n) is the number of divisors; p,( n) is the
Mobius function; A(n) is the Mangoldt function; a(n) is
the sum of divisors of the number n. Arithmetic func-
tions also include the integral part of a number, [x],
and the fractional part of a number, {x}. Arithmetic
functions giving the number of solutions of an equation
are also studied; for example, r (n) is the number of
integer solutions x and y of the equation x 2 +y2 = n in
the Goldbach problem; J (N) is the number of solutions
in prime numbers of the equation N=PI+P2+PJ'
Other arithmetic functions express the quantity of
numbers satisfying certain conditions; thus, for
instance, the function 7T(X) - the number of primes not
larger than x - describes the distribution of primes;
7T(X, q, I) gives the number of primes not larger than x
in the arithmetic progression p 1 (modq). The Che-
byshev functions also deal with properties of primes:

O(x) is the sum of the natural logarithms of the prime
numbers up to x, while "'(x) = ~ n~x A(n) (d. Che-
byshev function).
Algebraic number theory deals with generalizations
of the above arithmetic functions of a natural argu-
ment. Thus, for instance, in an algebraic field K of
degree n, Euler's function CP(U) - the number of resi-
due classes by the ideal U mutually prime with U - is
introduced for an integral ideal U.
Arithmetic functions appear and are employed in
studies on the properties of numbers. However, the
theory of arithmetic functions is also of independent
interest. The laws governing the variations of arithmetic
functions cannot usually be described by simple formu-
las, and the asymptotic behaviour in terms of numerical
functions is determined. Since many arithmetic func-
tions are not monotone, the study of their average
values is of great importance. An important class of
arithmetic functions is constituted by the multiplicative
arithmetic functions and the additive arithmetic func-
tions (cf. Multiplicative arithmetic function; Additive
arithmetic function). The problem of the distribution of
their values is studied in probabilistic number theory
[5].
References
[I] WINOGRADOW, I.M. [I.M. VINOGRADOV]: Elemente der
Zahlentheorie, Oldenburg, 1956 (translated from the Russian).
[2] VINOGRADOV, I.M.: The method of trigonometric sum, in th"
theory of numbers, Interscience, 1954 (translated from the Rus-
sian).
[3] HUA, L.-K.: 'Abschatzungen von Exponentialsummen und ihre
Anwendung in der Zahlentheorie', in Enzyklopaedie der
Mathematischen Wissenschaften mit Einschluss ihrer Anwen-
dungen. Vol. I, 1959, Heft 13, Teil I.
[4] CHANDRASEKHARAN, K.: Arithmeticalfunctions, Springer. 1970.
[5] KUBILYUS. I.: Probabilistic methods in the theory of numbers.
Amer. Math. Soc . 1964 (translated from the Russian).
N.I. Klimov
Editorial comments.
References
[A1] HARDY, G.H. and WRIGHT, E.M.: An introduction to the
theory of numbers, Oxford Univ. Press, 1979, Chapt.16.
AMS 1980 Subject Classification: 10A20
AluTuMETIc GENUS - A numerical invariant of
algebraic varieties (cf. Algebraic variety). For an arbi-
trary projective variety X (over a field k) all irreducible
components of which have dimension n, and which is
defined by a homogeneous ideal I in the ring
k[To, ... ,TN], the arithmetic genus Pa(X) is expressed
using the constant term cp(I, 0) of the Hilbert polyno-
mial cp(I, m) of I by the formula
Pa(X) == (-It(CP(I, 0)-1).
This classical definition is due to F. Severi [1]. In the
general case it is equivalent to the following definition:
ARITHMETIC GROUP
Pa(X) == (-It(X(X, (!) x)-I),
where n
X(X, (!) x) == ~(-IY dimk H;(X, (!) x)
;=0
is the Euler characteristic of the variety X with coeffi-
cients in the structure sheaf (!) x. In this form the defin-
ition of the arithmetic genus can be applied to any
complete algebraic variety, and this definition also
shows the invariance of Pa(X) relative to biregular map-
pings. If X is a non-singular connected variety, and
k = C is the field of complex numbers, then
n-I
Pa(X) == ~gn-I(X)'
;=0
where gk(X) is the dimension of the space of regular
differential k-forms on X. Such a definition for n = 1,2
was given by the school of Italian geometers. For
example, if n = 1, then Pa(X) is the genus of the curve
X; if n =2,
where q is the irregularity of the surface X, while Pg is
the geometric genus of X.
For any divisor D on a normal variety X, O. Zariski
(see [1]) defined the virtual arithmetic genus Pa(D) as the
constant term of the Hilbert polynomial of the coherent
sheaf (!) xeD) corresponding to D. If the divisors D and
D' are algebraically equivalent, one has
Pa(D) == Pa(D').
The arithmetic genus is a birational invariant in the
case of a field k of characteristic zero; in the general
case this has so far (1977) been proved for dimensions
n~3 only.
References
[I] BALDASSARRl, M.: Algebraic varieties, Springer, 1956.
[2] HIRZEBRUCH, F.: TopolOgical methods in algebraic geometry,
Springer, 1978 (translated from the German).] V D I h
. . 0 gac ev
AMS 1980 Subject Classification: 14A99
ARITHMETIC GROUP - A subgroup H of a linear
algebraic group G defined over the field Q of rational
numbers, that satisfies the following condition: There
exists a faithful rational representation p: G~GLn
defined over Q (d. Representation theory) such that
p(H) is commensurable with peG) n
GL(n, Z), where Z
is the ring of integers (two subgroups A and B of a
group C are called commensurable if A B is of finite n
index in A and in B). This condition is then also satis-
fied for any other faithful representation defined over
Q. More generally, an arithmetic group is a subgroup
of an algebraic group G, defined over a global field k,
that is commensurable with the group Go of O-points
of G, where 0 is the ring of integers of k. An arith-
metic group H n
GR is a discrete subgroup of GR.
If </>: G~Gl is a k-epimorphism of algebraic groups,
253
ARITHMETIC GROUP
then the image cj>(H) of any arithmetic group He G is
an arithmetic group in G) [1]. The name arithmetic
group is sometimes also given to an abstract group that
is isomorphic to an arithmetic subgroup of some alge-
braic group. Thus, if k is an algebraic number field, the
group Go~G~ , where G' is obtained from G by res-
tricting the field of definition from k to Q, is called an
arithmethic group. In the theory of Lie groups the
name arithmetic subgroups is also given to images of
arithmetic subgroups of the group of real points of G R
under the factorization of G R by compact normal sub-
groups.
References
[I) BOREL, A.: Ensembles fundamentaux pour les groups
arithmetiques et formes automorphes, Fac. Sc. Paris, 1967.
[2) BOREL, A. and HARISH-CHANDRA: 'Arithmetic subgroups of
algebraic groups', Ann. of Math. 75 (1962), 485-535.
[3) Arithmetic groups and discontinuous subgroups, Proc. Symp.
Pure Math., 9, Amer. Math. Soc., 1966.
VP. Platonov
Editorial comments. Useful additional references are
[A1] - [A3]. [A2] is an elementary introduction to the theory
of arithmetic groups.
Conjectures of A. Selberg and 1.1. Pyatetskii-Shapiro
roughly state that for most semi-simple Lie groups discrete
subgroups of finite co-volume are necessarily arithmetic.
G.A. Margulis settled this question completely and, in
particular, proved the conjectures in question. See Discrete
subgroup for more detail.
References
[Ai) BOREL, A.: 'Arithmetic properties of linear algebraic groups',
in Proc. Internat. Congress Mathematicians Stockholm, 1962,
Uppsala, 1963, pp. 10-22.
[A2) BOREL, A.: Introduction aux groupes arithmetiques, Hermann,
1969.
[A3) HUMPHREYS, J.E.: Arithmetic groups, Springer, 1980.
AMS 1980 Subject Classification: 20G15, 20G30
ARITHMETIC MEAN - The number a obtained by
dividing the sum of several numbers aI, ... an by
their number n:
a =
AMS 1980 Subject Classification: 10A99
ARITHMETIC PROGRESSION, arithmetic series of
the first order - A sequence of numbers in which each
term is obtained from the term immediately preceding
it by adding to the latter some fixed number d, which is
known as the difference of this pror,ression. Thus, each
arithmetic progression has the form
a. a+d, a+2d .... ,
in which the general term is
(1" = a + (n ~ l)d.
254
A characteristic property of an arithmetic progression is
an+1 +an -I
an = 2
If d>O, the progression is increasing; if d<O, it is
decreasing. The simplest example of an arithmetic pro-
gression is the series of natural numbers 1, 2, . . .. The
number of terms of an arithmetic progression can be
bounded or unbounded. If an arithmetic progression
consists of n terms, its sum can be calculated by the
formula:
Sn
Editorial comments. For results on prime numbers in
arithmetic progressions see Distribution of prime numbers.
AMS 1980 Subject Classification: 10A40
ARITHMETIC PROPORTION - An equation of the
form
a-b = c-d,
where a, b, c, d are given numbers. An arithmetic pro-
portion is also called a difference proportion.
AMS 1980 Subject Classification: 10A99
ARITHMETIC ROOT, arithmetical value of the n-th
root of a real number a;;'O - A non-negative number the
n-th power of which is equal to a. In considering the
two real values of the n-th root, where n is an even
number, of a non-negative number, one speaks of the
algebraic values of the root in the domain of the real
numbers; if, on the other hand, one is considering all n
values of an n-th root, then one speaks of the values of
the root in the domain of the complex numbers.
AMS 1980 Subject Classification: 26A09
ARITHMETIC SERIES of order m - The sequence of
values of a polynomial of degree m:
p(x) = ao+a\x+ ... +amx m ,
assumed by the polynomial when the variable x takes
successive integral non-negative values x = 0, 1, .... If
m=l, i.e. p(x)=aO+alx, one obtains an arithmetic
progression with initial term ao and difference a I. If
p(x)=x 2 or p(x)=x 3 , one obtains sequences of
squares or cubes of integers. i.e. special cases of arith-
metic series of the second and third orders. If a first
difference series is created, constituted by the differ-
ences between successive terms of an arithmetic series,
then a series of differences of the first difference series
(second differences) is written. and then the second
differences are used to form third differences, etc .. one
finds. at the moth stage, that all moth differences are
equal. Conversely. if the m-th differences of a numeri-
cal sequence are all equal. the sequence is an arithmetic
series of order m. Using this property. it is possible to
 ARITH\-IFIICAL A 'vTRAG ES. SliMMATION METHOD OJ

construct arithmetic series of different orders from their ARITHMETIC SPACE, number space, coordinate
differences. For example, the sequence I, I, I, ... , space, real n-space - A Cartesian power R" of the set of
may be regarded as the first differences of the series of real numbers R having the structure of a linear topo-
natural numbers I, 2, 3, ... ; as the second differences logical space. The addition operation is here defined by
of the series of triangular numbers I, 3, 6, 10, ... ; as the the formula:
third differences of the sequence of tetrahedral numbers
1,4, 10,20, ... , etc. These numbers are so called
because triangular numbers represent numbers while multiplication by a number AER is defined by
arranged in the form of a triangle (Fig. I), while the formula
tetrahedral numbers represent numbers arranged in the A(X" ... ,xn ) = (Ax"". , Axn ).
form of tetrahedra (pyramids) (Fig. 2).
The topology in Rn is the topology of the direct pro-
duct of n copies of R; its base is formed by open n-
a
1
4~ 10
dimensional parallelepipeda:
I = {(x" ... ,xn)ERn: aj<x,<bj, i=l, ... on},
Fig. 1. Fig. 2.
where the numbers a j, . . . ,an and b 1, . . . ,bn are
Triangular numbers are expressed by the formula given.
[n (n + 1)1/ 2, while tetrahedral numbers are given by The real n-space Rn is also a normed space with
the formula respect to the norm
n(n + l)(n +2) n = 1, 2, ....
II x I = v'XT + ... + x~,
6
where x = (x 1, . . . , Xn) E Rn , and is a Euclidean space
A generalization of triangular numbers is constituted with respect to the scalar product
by k-gonal or figurate numbers, which played an impor-
tant role in the development of arithmetic in its various <x,y> = x,y"
i ---::: 1
stages.
wherex=(Xj, ... ,xn),Y=(yj, ... ,Yn)ERn.
I. V Dolgachev
a~-a
~-~~ AMS 1980 Subject Classification: 15A03, b1 N20,
.-~-4t
9
468XX

Fig. 3. Fig. 4.
ARITHMETIC TRIANGLE - A synonym for Pascal
triangle.
k-gonal numbers are given by the formula:
_ n(n-1) _ AMS 1980 Subject Classification: 10A99
n +(k 2) 2 ' n -1, 2, ....

They form an arithmetic series of the second order, ARITHMETICAL AVERAGES, SUMMATION
with one as their first term, k as their second term and METHOD OF - One of the methods for summing series
k - 2 as their second differences. If k = 3, triangular and sequences. The series
numbers are obtained; if k = 4, one obtains squares
(n 2); if k = 5, one obtains pentagonal numbers
(3n 2 - n) / 2, etc. These appellations will become clear
is summable by the method of arithmetical averages to
from Figs. 3 and 4, in which the number of beads
the sum s if
arranged in the form of a square or a pentagon is s 0 + ... +s n
Ii
expressed by the respective square or pentagonal n~n;, n +1 = s,
number. Figurate numbers satisfy the following
theorem, proposed by P. Fermat and first proved by where Sn = ~~ = 1Uk In this case, one also says that the
A.L. Cauchy: Any natural number can be represented sequence {sn} is summable by the method of arithmeti-
as a sum of not more than k k-gonal numbers. cal averages to the limit s. The summation method of
arithmetical averages is also called the Cesaro summa-
References
[I) WAERDEN, B.L. VAN DER: Algebra, 1-2, Springer. 1967-1971
tion method of the first order (cf. Cesaro summation
(translated from the German). methods). The summation method of arithmetical aver-
[2) ARNOL'D, V.1.: Theoretical arithmetic, Moscow, 1939 (in Rus- ages is completely regular (see Regular summation
sian).
BSE-2 methods) and translative (see Translativity of a summa-
AMS 1980 Subject Classification: 10A40 tion method).

255
ARITHMETICAL AVERAGES, SUMMATION METHOD OF
References
[1] HARDy, G.H.: Divergent series, Clarendon, 1949.
1.1. Volkov
Editorial comments. Instead of 'arithmetical averages'
the term 'arithmetical means' is sometimes used, cf. [A 1],
and instead of 'summation' one also uses 'summability':
summability method.
References
[A1] COOKE, R.G.: Infinite matrices and sequence spaces, Macmil-
lan, 1950.
AMS 1980 Subject Classification: 40G05
AluTHMETIZATION - A method used in mathemat-
ical logic for replacing a reasoning on the expressions
of some logico-mathematical language by reasonings on
natural numbers. For this purpose the replacement is
constructed by some sufficiently simple one-to-one
mapping of the set of all words (in the alphabet of the
language under consideration) into the natural number
sequence. The image of a word is called its number.
Relations between and operations defined on words are
transformed by this mapping into relations between
and operations on natural numbers. The requirement of
a 'sufficiently simple' mapping leads to the fact that
some basic relations (such as the relation of imbedding
of one word into another, etc.) and some operations
(like the operation of concatenation of words, etc.) are
transformed into relations and operations having a sim-
ple algorithmic nature (e.g. are primitive recursive). In
particular, if among the expressions of the language
under consideration there exist programs for some fam-
ily of computable functions (d. Computable function),
arithmetization naturally leads to an enumeration of
this family (in which for the number of each function is
taken the number of its program).
The first arithmetization was constructed by K.
G6del [1] for the proof of the incompletness of formal
arithmetic (d. GOdel incompleteness theorem). More
precisely, G6del put the letters of the alphabet in
correspondence with some, pairwise different, natural
numbers and attached to the word 1') . Tn the
number i' ... p~, where ti is the number correspond-
ing to the letter Ti' and Pi is the i-th prime number in
the natural number sequence. Such an enumeration is
called a Codel enumeration. In a broad sense every
enumeration of words arising from arithmetization is
called a Gbdel enumeration and the number
corresponding to a word is called its Codel number.
In 1936 A. Church used arithmetization to obtain the
first example of an unsolvable algorithmic problem of
arithmetic.
The term 'arithmetization' (in the phrase 'arithmeti-
zation of analysis') is also used in the literature on the
foundations of mathematics for the denotation of the
256
creation of the theory of real numbers in the nineteenth
century using set-theoretic constructions, starting from
the natural numbers.
References
[I] GOOEL, K.: 'Ueber formal unentscheidbare Satze der Principia
Mathematica und verwandter Systeme 1', Monatsh. Math. Phy-
sik 38 (1931), 178-198.
[2) CHURCH, A.: 'An unsolvable problem of elementary number
theory', Amer. J. Math. 58 (1936), 345-363.
[3) KLEENE, S.c.: Introduction to metamathematics, North-Holland,
1951.
V.A. Uspenskil
Editorial comments.
[A1) SMORINSKY, c.: 'The incompleteness theorems', in J. Barwise
(ed.): Handbook of mathematical logic, North-Holland, 1977,
pp. 821-866.
AMS 1980 Subject Classification: 03F40
ARRANGEMENT with repetitions of n out of m ele-
ments - A finite sequence a= (ail' ... , aiJ of elements
from some set A ={a), ... ,am}. If all the term of a
are distinct, then a is called an arrangement without
repetitions. The number of arrangements of n out of m
elements with repetitions is m n , and that without
repetitions in (m)n = m(m - 1) ... (m - n - 1).
An arrangement can be regarded as a function </>
given on Zn = { 1, ... ,n} and taking values in A:
</>(k)=ai k=l, ... ,no The elements of A are usually
k '
called cells (or urns), while those of Zn are called parti-
cles (or balls); </> defines the filling of the various cells
by the various particles. If one speaks of indistinguish-
able particles or cells, it is understood that classes of
arrangements are being considered. Thus, if all the par-
ticles are indistinguishable, then two arrangements
defined by functions </> and 1/;, respectively, belong to
the same class if and only if there is a permutation a of
Zn such that </>(a(k=I/;(k) for all k EZn. In this case,
the number of such classes, or, as it is called, the
number of arrangements filling m different cells by n
indistinguishable particles, is the number of combina-
tions of n out of m elements, with repetition allowed
(d. Combination).
If one says that all the cells are indistinguishable,
then one means that the arrangements are put into the
classes in such a way that two arrangements, defined
by functions </> and I/; respectively, belong to the same
class if there exists a permutation l' of A for which
T(</>(k=I/;(k) for all k EZn- In this case the number of
arrangements of n different particles in m indistinguish-
able cells, that IS, the number of classes, IS
~;=)S(n,k), where S(n,k) is the Stirling number of
the second kind:
S(n,k) = S(n-l,k-I)+kS(n-l,k), S(O, 0) = I,
S(n,k) = fork>n andk=O, n>O.

If one does not distinguish either particles or cells, then
an arrangement of n identical particles in m identical
cells is obtained; the number of such arrangements is
L.~=IPn(k), where Pn(k) is the number of additive par-
titions of n into k natural numbers. One can also con-
sider other partitions of arrangements into classes, for
example when the above-mentioned permutations (J and
T are taken from subgroups of the symmetric groups of
degrees n and m, respectively (this and other generaliza-
tions can be found in [I], [2]). Synonyms of 'arrange-
ment' are the terms 'n-permutation', and 'ordered n-
sample of a population'.
References
[I] SACHKOV, A.V.: Combinatorial methods of discrete mathematics,
Moscow, 1977 (in Russian).
VM. Mikheev
Editorial comments.
References
[A 1] COMTET, L.: Advanced combinatorics, Reidel, 1977.
AMS 1980 Subject Classification: 05A05, 05A 17
ARTINIAN GROUP, group with the minimum condi-
tion for subgroups, group with the descending chain condi-
tion - A group in which any decreasing chain of distinct
subgroups terminates after a finite number. Artinian
groups are periodic, and the question of their structure
hinges on Schmidt's problem on infinite groups with
finite proper subgroups [3] and the minimality problem:
Is an Artinian group a finite extension of an Abelian
group? Both these problems have been solved for
locally solvable groups [I] and locally finite groups [3],
[4].
References
[I] CHERNIKHOV, S.N.: 'Infinite locally solvable groups', Mat. Sb.
7 (49), no. I (1940), 35-64 (in Russian).
[2] CHERNIKHOV, S.N.: 'The finiteness condition in general group
theory', Uspekhi Mat. Nauk 14, no. 5 (1959), 45-96 (in Rus-
sian).
[3] KARGAPOLOV, M.I.: 'On a problem of O.Yu. Schmidt', Sibirsk.
Mat. Zh. 4, no. I (1963),232-235 (in Russian).
[4] SHUNKOV, V.P.: 'On the minimality property for locally finite
groups', Algebra and Logic 9, no. 2 (1970), 137-151
(Algebra i Logika 9, no. 2 (1970), 220-248)
VP. Shunkov
Editorial comments. Schmidt's problem actually states:
Under what conditions does an infinite group have proper
infinite subgroups?
References
[A1] KEGEL, O.H. and WEHRFRITZ, B.V.: Locally finite groups,
North-Holland, 1973.
AMS 1980 Subject Classification: 20FXX, 20E07,
20E28
ARTINIAN MODULE - A module that satisfies the
decreasing chain condition for submodules. The class of
Artinian modules is closed with respect to passing to
ARTINIAN RING
submodules, quotient modules, finite direct sums and
extensions. Extension in this context means that if the
modules B and A / Bare Artinian, then so is A. Each
Artinian module can be decomposed into a direct sum
of sub modules which are no longer decomposable into
a direct sum. A module has a composition series if and
only if it is both Artinian and Noetherian. See also
Artinian ring.
L.A. Skornyakov
Editorial comments.
References
[A1] FAITH, c.: Algebra: rings, modules and categories, 1,
Springer, 1973.
[A2] FAITH, c.: Algebra II: ring theory, Springer, 1976.
AMS 1980 Subject Classification: 13E1 0
ARTINIAN RING, right Artinian ring - A ring that
satisfies the minimum condition for right ideals, i.e. a
ring in which any non-empty set M of right ideals that
is partially ordered by inclusion has a minimal element
[1] - a right ideal from M that does not strictly con-
tain right ideals from M. In other words, an Artinian
ring is a ring which is a right Artinian module over
itself. A ring A is an Artinian ring if and only if it
satisfies the decreasing chain condition for right ideals,
i.e. for any decreasing sequence of right ideals
Bid B 2 d . .. there exists a natural number m such
that Bm = Bm + I = .... The definition of a left Artinian
ring is similar.
Each associative Artinian ring with a unit is right
Noetherian (cf. Noetherian ring). Every finite-
dimensional algebra over a field is an Artinian ring.
The properties of Artinian rings in the class of alterna-
tive rings, and particularly in the class of associative
rings, have been studied most fully (cf. Alternative rings
and algebras; Associative rings and algebras). The
Jacobson radical of an associative Artinian ring is nil-
potent and contains every one-sided nil-ideal. A ring A
is a simple associative Artinian ring if and only if it is
isomorphic with the ring of all matrices of some finite
order over some associative skew-field. In the class of
alternative rings any simple Artinian ring is either asso-
ciative, or else is a Cayley - Dickson algebra over its
centre, which is then a field. The structure of associa-
tive Artinian rings with zero Jacobson radical has been
described (cf. Semi-simple ring). There exists a variant
of this theorem in the case of alternative rings. For
associative rings with non-zero Jacobson radical a
fairly-advanced structure theory has been developed [1],
[2]. Several classes of Artinian rings - quasi-Frobenius
rings, uniserial rings, balanced rings (cf. Quasi-
Frobenius ring; Uniserial ring; Balanced ring) - are
being intensively studied.
References
[I] ARTIN, E., NESBITT, C.l. and THRALL, R.M.: Rings with
257
ARTINIAN RING

minimum condition, Univ. Michigan Press, Ann Arbor, 1946. References

(2) JACOBSON, N.: Structure of rings, Amer. Math. Soc., 1956. (1) ARzELA, C: Rend. Accad. Sci. Bologna 9, no. 2 (1905), 100-
(3) Itogi Nauk. Algebra Topol. Geom. 1965 (1967), 133-180. 107.
(4) 1togi Nauk. Algebra Topol. Geom. 1968 (1970),9-56. (2) HAHN, H.: Theorie der reel/en Funktionen, 1, Springer, 1921.
K.A. Zhevlakov B.1. Golubov
Editorial comments. AMS 1980 Subject Classification: 26830
References
[A 1) FAITH, C: Algebra: rings, modules and categories, 1, ASSERTION, proposition, statement, (assertive) sen-
Springer, 1973. tence - A narrative communication which by virtue of
[A2] FAITH, c.:Algebra II: ring theory, Springer, 1976.
its meaning is true or false. In a narrower sense, the
AMS 1980 Subject Classification: 13E1 0 term assertion in mathematical logic is understood to
be a closed formula in a logico-mathematical language
ARZELA - ASCOLI THEOREM - The name of a which by virtue of the semantics of the language can be
number of theorems that specify the conditions for the classified as true or false.
limit of a sequence of continuous functions to be a con- Thus, in axiomatic set theory, various mathematical
tinuous function. One such condition is the quasi- statements, for example, the axiom of choice or the con-
uniform convergence of the sequence. tinuum hypothesis, are written in the form of formulas
References that for general semantic reasons can be interpreted as
[1) ARZELA, C: Mem. Accad. Sci. Bologna (5) 5 (1893),225-244. expressing the contents of the statement. It is by no
[2) ASCOLI, G.: Rend. Accad. Lincei 18 (1883),521-586.
means necessary that there exists a method of recogniz-
P.s. Aleksandrov ing which assertions in the language are true and which
Editorial comments. false. Moreover, the semantics need not actually be
References sufficiently developed, or may present fundamental dif-
[A 1) DUNFORD, N. and SCHWARTZ, IT.: Linear operators. General ficulties in the solution of the question of whether cer-
theory, 1, Interscience, 1958. tain assertions are true. The un solvability of assertions
within the framework of a theory is explained by the
ARZELA VARIATION - A numerical characteristic
formalization method (for examples, see Axiomatic set
of a function of several variables, which may be
theory).
regarded as the multi-dimensional analogue of the vari- A. G. Dragalin
ation of a function III one unknown. Let AMS 1980 Subject Classification: 00A25, 03A05
f(x)=f(x], ... ,xn) be a rea)-valued function given on
an n-dimensional parallelepipedon ASSOCIATED FUNCTION of a complex variable - A
D n =[a1,bJlX'" X[an,b n ], n=2,3, ... , and let Gn function which is obtained in some manner from a
be the class of all continuous vector functions given function f (z) with the aid of some fixed function
x(t)=(X1(t), ... ,xn(t, O.:;;;t.:;;;l, such that each of the F(z). For example, if
functions Xk is non-decreasing on [0, 1], and with
= :L ak zk
00

xk(O)=ako xk(l)=bk , k = 1, ... ,no Then f(z)

k =0
m
is an entire function and if
A(j,Dn) = _sup S~? ~ If(x(t,))-fCx(t,-d) I,
x(t)""G/I =1
= :L bkz k
S 00

F(z)
where II={0=t o <t1<'" <tm =1} is an arbitrary k =0

system of points in [0, 1]. This definition for the case is a fixed entire function with bk*O, k~O, then
n =2 was proposed by C. Arzela [I] (see also [2], p. oc ak
543). If A(j, Dn)<oo, one says thatfhas bounded (fin- y(z) = :L-Z-(k+I)
k =0 bk
ite) Arzela variation on D n , and the class of all such
functions is denoted by A (Dn). For a function is a function which is associated to f (z) by means of
f(x 1, . . . ,xn) to helong to the class A (Dn) it is neces- the function F(z); it is assumed that the series con-
sary and sufficient that there exists a decomposition verges in some neighbourhood I z I > R. The function
f= fl -h, where fl and h are finite non-decreasing f (z) is then represented in terms of y(z) by the formula
functions on Dn. A function f is called non-decreasing f(z) = 1.
-2 f y(t)F(zt)dt.
on Dn if 'TTl It I =R,>R
f(x;, ... ,x;;) ~ f(Xj, . ... x;,),
In particular, if ak
for ak ':;;;x~ <x~ .:;;;bk (k = I, ... ,IJ). The class A (Dn) f(z) = :L -zk
00

contains the class of functions of bounded Hardy varia- k=(j k!

tion on Dn. is an entire function of exponential type and F(z)=e z ,

251-\

then
=
'Xl
y(z) ~akZ-(k+I)
k ~O
is the Borel-associated function of f (z) (cf. Borel
transform).
A.F. Leont'ev
AMS 1980 Subject Classification: 30BXX, 44AXX
ASSOCIATIVE CALCULUS - A name given to cal-
culi (cf. Calculus) of a very definite type, especially suit-
able for specifying finitely-presented associative systems
(semi-groups, cf. Semi-group). The term 'associative
calculus' was introduced by A.A. Markov, who also
developed the theory of associative calculi [2].
Any associative calculus m: is defined by specifying
some alphabet A and a finite list a of relations in A -
pairs of words (cf. Word) in this alphabet. The words
constituting a relation are often called its parts - left
and right. A a-permissible operation on words in A is
any substitution of one part of any relation belonging
to a, where it occurs imbedded in a word in m: (cf.
Imbedded word), by the other part of the same relation.
An associative calculus m: is the performance of a
sequence of a-permissible operations, starting from any
word P in A. One says that all words Q obtained in
this way (including the initial word P itself) are
equivalent to P in the associative calculus m: (in sym-
bolic notation, m:: PrrQ). The relation rr for any associ-
ative calculus m: is reflexive, symmetric and transitive.
In addition, it follows from m:: PrrQ and m:: RrrS that
m:: PRrrQS. This makes it possible, in a natural
manner, to assign to each associative calculus m: a
finitely-presented associative system Kill:, whose ele-
ments are equivalence classes of words in A and whose
multiplication operation is induced by the operation of
concatenation of words in A. The associative system
Kill: thus constructed will have a unit (the element
represented by the empty word); the elements of K &
represented by the letters of the alphabet A will consti-
tute a system of generating elements for Kill:, while the
list of relations a represents a complete system of rela-
tions between these generating elements of Kill: in the
sense that the elements of Kill: represented by two
words P and Q are identical in Kill: if and only if P and
Q are equivalent in m:. Thus, the recognition of the
identity of elements in Kill: is reduced to the recognition
of equivalence of words in m:. Hence the importance of
the study of the solvability of the algorithmic problem
of recognition of equivalent words in an arbitrary asso-
ciative calculus. This problem was first formulated by
A. Thue [1]: Construct, for an arbitrary associative cal-
culus m:, an algorithm which, for any pair of words in
the alphabet of this associative calculus, determines, in
a finite number of steps, whether or not the words con-
ASSOCIA TIVE CALCULUS
stituting the pair are equivalent in m:. In the algebraic
interpretation, this problem is the identity problem for
the associative system Kill:. Thue succeeded in solving
the problem in a few special cases only. In the modern
interpretation of the problem (1930's and later) the
algorithm of this problem is sought in some exact
meaning of the word, such as a partial recursive
function, a Turing machine or a normal algorithm. If the
problem is thus modernized, the problem of finding a
specific associative calculus for which the desired algo-
rithm does not exist becomes meaningful. Markov [3]
and E.L. Post [4], working independently of each other,
constructed unsolvable associative calculi, i.e. associative
calculi with an unsolvable algorithmic problem of
recognition of word equivalence. These results yield a
negative solution of Thue's problem in its modern
form. However, if one accepts the Church thesis or any
other equivalent assumption regarding the adequacy of
recursive functions for rendering the concept of an
algorithm more precise, one comes to the inevitable
conclusion that Thue's problem in its original form as
well has a negative solution for some concrete associa-
tive calculi.
The original examples constructed by Markov and
Post were extremely complex. Simpler unsolvable asso-
ciative calculi were subsequently presented. They
include, for example, an associative calculus with seven
very simple relations [5], and one with only three such
relations [6]. Thue's problem has been almost fully
solved for the case of an associative calculus with one
relation [7].
In a natural way one defines an isomorphism of one
associative calculus into another and onto another [2].
Of special interest from the algebraic point of view is
the study of those properties of associative calculi that
are invariant under isomorphisms; they are the
properties of abstract associative systems. Markov [2],
who based himself on his own studies on the recogni-
tion of word equivalence in an associative calculus,
obtained a very general result, which yields a negative
solution to practically all algorithmic problems con-
cerning the basic classifications of associative calculi
studied at that time. He showed, in particular, that if I
is an isomorphism-invariant property of an associative
calculus, if there is a unique associative calculus with
the property I and if there exists also an associative
calculus not included in any associative calculus with
the property I, then the algorithmic problem of recog-
nition of an associative calculus with the property I
among other associative calculi is unsolvable for any
alphabet with more than three letters. It immediately
follows that algorithmic problems of recognition of
uniqueness of associative calculi, finiteness of an associ-
ative calculus, the semi-group character of an associa-
259
ASSOCIA T1VE CALCULUS
tive calculus, inclusion in group-associative calculi, and
isomorphism of pairs of associative calculi, are unsolv-
able for any alphabet with more than three letters. It
was subsequently proved [8] that the set of pairs of iso-
morphic associative calculi is recursively enumerable;
the method used in this connection also makes it possi-
ble to establish the recursive enumerability of the
number of invariant properties of associative calculi.
References
[lA) TRUE, A.: 'Probleme tiber Veriinderungen von Zeichenreihen
nach gegebener Regeln', Kra. Vidensk. Selsk. Skrifter. 1. Mat.
Nat. KI., no. 10 (1914).
[lB) TRUE, A.: 'Probleme tiber Veriinderungen von Zeichenreihen
nach gegebener Regeln', in Selected Math. Papers, Univ. For-
laget, Oslo, 1977, pp. 493-524.
[2) MARKov, A.A.: Theory of algorithms, Israel Progr. Sci. Transl.,
Jerusalem, 1961 (translated from the Russian). Also: Trudy
Mat. Inst. Steklov. 42 (1954).
[3) MARKov, A.A.: 'On the impossibility of certain algorithms in
the theory of associative systems', c.R. Acad. Sci. URSS 55
(1948),533-586. (Dokl. Akad. Nauk SSSR 55, no. 7 (1947),
587-590)
[4) POST, E.L.: 'Recursive unsolvability of a problem of Thue', J.
Symbolic Logic 12, no. I (1947), I-II.
[5) TSEITIN, G.S.: 'Associative calculi with undecidable
equivalence problems', Dokl. Akad. Nauk SSSR 107, no. 3
(1956), 370-371 (in Russian).
[6) MATIYASEVICH, YU.V.: 'Simple examples of undecidable asso-
ciative calculi', Soviet Math. Dokl. 8, no. 2 (1967), 555-557.
(Dokl. Akad. Nauk SSSR 173, no. 6 (1967), 1264-1266)
[7) ADJAN, S.I. [S.I. ADYAN): 'Defining relations and algorithmic
problems for groups and semigroups', Proc. Steklov Inst. Math.
85 (1967). (Trudy Mat. Inst. Steklov. 85 (1966
[8) NAGORNYi, N.M.: 'On the search of isomorphisms of associa-
tive calculi', Z. Math. Logik Grundl. Math. 6 (1960),319-324
(in Russian). German abstract. NM N
. . agornYI
Editorial comments. The term 'associative calculus'
seems to be restricted to the Russian literature; in the West
the term 'Thue system' is often used. For an account of
Thue systems and related combinatorial decision problems,
see [A1], Chapt. 6, or [A2].
References
[A 1) DAVIS, M.: Computability and unsolvability, McGraw-Hili,
1958.
[A2) DAVIS, M.: 'Unsolvable problems', in J. Barwise (ed.): Hand-
book of mathematical logic, North-Holland, 1977, pp. 567-
594.
AMS 1980 Subject Classification: 03003
ASSOCIATIVE RINGS AND ALGEBRAS - Rings
and algebras with an associative multiplication, i.e. sets
with two binary operations, addition + and multiplica-
tion " that are Abelian groups with respect to addition
and semi-groups with respect to multiplication, and in
which the multiplication is distributive (from the left
and from the right) with respect to the addition. More-
over, an associative algebra should be a. vector space
over a fixed field F, and the multiplication in it is com-
patible with multiplication by field elements in the
sense of the condition a(ab)=(aa)b=a(ab) for all aEF
260
and all a, b from the algebra.
The first examples of associative rings and associative
algebras were number rings and fields (the field of
complex numbers and its subrings), polynomial alge-
bras, matrix algebras over fields, and function fields.
The theory of associative rings and algebras became an
independent part of algebra at the beginning of the 20-
th century. This theory has many contact points with
numerous fields of mathematics, especially so with alge-
braic geometry and algebraic number theory (commuta-
tive rings), functional analysis (commutative normed
rings, operator rings and function rings), and topology
(rings of continuous functions on topological spaces).
The theory of fields, the theory of commutative rings
(cf. Field; Commutative ring, see also Commutative
algebra) and the theory of representations of associative
algebras have become independent branches of the
theory of associative rings and algebras. The theory of
topological rings and skew-fields forms part of topologi-
cal algebra.
The classical part of the theory of associative rings
and algebras is formed by the theory of finite-
dimensional associative algebras [2]. The principal
results of this theory are: A finite-dimensional simple
associative algebra (i.e. one without proper ideals) over
a field F is a complete matrix algebra over some skew-
field that is finite-dimensional over F (Wedderburn's
theorem); a finite-dimensional associative algebra over a
field of characteristic zero (and, even more generally -
- a separable finite-dimensional associative algebra) is the
direct sum (as linear spaces) of its radical I (i.e. its
maximal nilpotent ideal) and some semi-simple (i.e.
with a radical zero) sub algebra S, any two complemen-
tary semi-simple sub algebras Sand S 1 being conjugate
(cf. Wedderburn- Mal'tsev theorem).
One of the most important classes of associative alge-
bras are skew-fields (cf. Skew-field, i.e. associative rings
in which the equations ax = band ya = b are solvable
for all a, b in the ring, a=r!=O). Skew-fields that are alge-
bras over a certain field are called division algebras (cf.
Division algebra). The theory of finite-dimensional divi-
sion algebras forms a classical part of the theory of
fields. All finite-dimensional associative division alge-
bras over the field of real numbers have been
described: these are the field of real numbers, the field
of complex numbers and the skew-field of quaternions
(cf. Frobenius theorem). All finite skew-fields are com-
mutative (Wedderburn's theorem on skew-fields). A
Galois theory of skew-fields has been constructed [5].
The key concepts in the structural theory of associa-
tive rings are the concepts of the Jacobson radical.
semi-simplicity and primitivity. An associative ring is
called semi-simple (in the sense of Jacobson) if its
Jacobson radical is zero. A ring is called primitive

(from the right) if it has an irreducible exact right
module. All semi-simple associative rings are a sub-
direct sum of primitive rings. Any primitive associative
ring R is a dense ring of linear transformations of some
vector space V over a skew-field (Jacobson's density
theorem); here the meaning of density is that for all
linearly independent elements v I, . . . ,Vk from V and
all elements U I, . . . ,Uk from V there exists a transfor-
mation rER such that Vjr=Uj for l';;;;i.;;;;k. The general
theory of radicals occupies an important place in the
structural theory of rings (cf. Radical of rings and alge-
bras).
The theory of Artinian rings (from the right, cf.
Artinian ring), i.e. rings with the descending chain con-
dition (minimum condition) for right ideals, forms a
classical part of the theory of associative rings. The
main result of this theory is that an associative ring is a
semi-simple Artinian ring if and only if it is a direct
sum of a finite number of complete matrix rings over
skew-fields (the Wedderburn- Artin theorem).
The concept of the (classical) ring of fractions is
important in the structural theory of associative rings.
A ring Q(R) is called the (right) ring of fractions of a
subring R of it if in Q(R) all the regular elements (i.e.
not divisors of zero) of R are invertible and if any ele-
ment of Q(R) has the form ab- I , where a, bER. An
associative ring has a ring of fractions if and only if for
all elements a, b ER, where b is regular, there exist ele-
ments aJ, b l ER such that ab l =bal and b l is regular
(Ore's theorem). A ring R has a semi-simple Artinian
ring of fractions if and only if it is semi-primary (i.e.
/2=1=0 for all non-zero ideals / c;;;,R), satisfies the
minimum condition for right annihilating ideals of the
form
r(S) = {xER: Sx=O},
where S is a subset of R, and contains no infinite direct
sums of right ideals (Goldie's theorem). In addition to
the classical rings of fractions, rings of fractions having
other meanings are also studied - mainly maximal or
complete rings of fractions [8].
Special stress is laid on the study of free associative
algebras (cf. Free associative algebra). Let F be a field
and let X be a set. The free associative algebra F <X >
over F with a unit element and with base X is the alge-
bra of non-commutative polynomials with free terms in
the set of variables X and with coefficients from F. The
algebra F < X> is characterized by the fact that it is is
generated by X as an algebra with a unit element and
any mapping of X into an associative algebra R with a
unit element can be extended (even uniquely) to a
homomorphism from F<X> into R. A free
associative algebra is a ring with free ideals, i.e. the
right (left) ideals of the ring F <X > are free right
(left) F <X>-modules, and all bases of the free
ASSOCIATIVE RINGS AND ALGEBRAS
finitely-generated F<X>-modules contain the same
number of elements (Cohn's theorem). Other examples
of rings with free ideals are group algebras of free
groups, and free products of associative division alge-
bras. The free associative algebra F <X > is also a
unique factorization domain: Any non-invertible ele-
ment aEF<X> has a representation a=pI .. 'pn,
where pj are irreducible elements, and this representa-
tion is unique except for the order of the sequence and
similarities (two element a and b of a ring R are called
similar if R / aR and R / bR are isomorphic as right
R-modules). The centralizer of each non-scalar element
of the algebra F < X> is isomorphic to the algebra of
polynomials F[t] in one variable t (Bergman's theorem).
Group algebras and PI-algebras are important classes
of associative algebras (cf. Group algebra; PI-algebra).
The theory of varieties of rings is in the process of
being developed.
The importance of ring theory in mathematics has
increased with the development of homological algebra.
Many known classes of rings may be characterized in
terms of properties of the category of modules (cf.
Modules, category of) over these rings. For instance, a
ring R is a semi-simple Artinian ring if and only if all
right (left) modules over R are projective (injective). A
ring R is regular (in the sense of von Neumann) if and
only if all right (left) modules over R are flat. See also
Regular ring; Homological classification of rings;
Quasi-Frobenius ring.
References
[1] WAERDEN, B.L. VAN DER: Algebra, 1-2, Springer, 1967-1971
(translated from the German).
[2] ALBERT, A.A.: Structure of algebras, Amer. Math. Soc., 1939.
[3] ARTIN, E., NESBITT, C.l. and THRALL, R.M.: Rings with
minimum condition, Univ. Michigan Press, Ann Arbor, 1946.
[4] JACOBSON, N.: The theory of rings, Amer. Math. Soc., 1943.
[5] JACOBSON, N.: Structure of rings, Amer. Math. Soc., 1956.
[6] KUROSH, A.G.: Lectures on general algebra, Chelsea, 1963
(translated from the Russian).
[7] HERSTEIN, I.: Noncommutative rings, Math. Assoc. Amer., 1968.
[8] LAMBEK, J.: Lectures on rings and modules, Blaisdell, 1966.
[9] CURTIS, C.W. and REINER, I.: Representation theory offinite
groups and associative algebras, Interscience, 1962.
[10] PONTRYAGIN, L.S.: Topological groups, Princeton Univ. Press,
1958 (translated from the Russian).
[II] COHN, P.M.: Free rings and their relations, Academic Press,
1971.
[12] ZARISKI, O. and SAMUEL, P.: Commutative algebra, 1-2, van
Nostrand, 1960. (reprinted: Springer, 1975).
[13] CARTAN, H. and ElLENBERG, S.: Homological algebra, Prince-
ton Univ. Press, 1956.
[14] Itogi Nauk. Algebra Topol. 1962 (1963), 59-79.
[15] Itogi Nauk. Algebra Topol. Geom. 1965 (1967), 133-180.
[16] Itogi Nauk. Algebra Topol. Geom. 1968 (1970),9-56.
[17] BOKUT', L.A., KUZ'MIN, E.N. and SHIRSHOV, A.I.: Rings, 1-3,
Novosibirsk, 1973 (in Russian).
[18] DIVINSKY, N.J.: Rings and radicals, Allen and Unwin, 1965.
[19] PASSMAN, D.S.: Infinite group rings, New York, 1971.
[20] PROCESI, C.: Rings with polynomial identities, New York, 1973.
L.A. Bokut'
261
ASSOCIATIVE RINGS AND ALGEBRAS

Editorial comments. from the point A is

References
[A1] FAITH, c.: Algebra I: Rings, modules and categories, I = ; R sin2 t.
Springer, 1973.
[A2] FAITH, c.: Algebra If: Ring theory, Springer, 1976.
The length of the entire curve is 6R. The radius of cur-
[A3] GILMORE, R.: Multiplicative ideal theory, M. Dekker, 1972. vature is
[A4] McDoNALD, B.R.: Finite rings with identity, M. Dekker, 1974. 'k = 23 R'sm 2t
AMS 1980 Subject Classification: 16-XX The area bounded by the curve is

ASSOCIATIVITY, law of associativity - A property of s = .l'ITR2

8 .
an algebraic operation. For the addition and multipli-
The astroid is the envelope of a family of segments of
cation of numbers associativity is expressed by the fol-
constant length, the ends of which are located on two
lowing identities:
mutually perpendicular straight lines. This property of
(a+b)+c = a+(b+c) and (ab)c = a (bc). the astroid is connected with one of its generalizations
A binary algebraic operation * is associative (or, - the so-called oblique astrOid, which is the envelope of
which is the same thing, satisfies the law of associativity) the segments of constant length with their ends located
if the identity on two straight lines intersecting at an arbitrary angle.
(X*y)*z = x*(y*z)
References
is valid in the given algebraic system. In a similar [I] SAVELOV, A.A.: Plane curves, Moscow, 1960 (in Russian).
manner, associativity of an n-ary operation w is defined D.D. Sokolov
by the identities Editorial comments.
References
= XI .. Xj(Xj+1 . . . Xj+nw)Xj+n+1 . X2n-IW, [A1] LAWRENCE, J.D.: A catalog of special plane curves, Dover,
reprint, 1972.
for all i= 1, ... ,no [A2] LOCKWOOD, E.A.: A book of curves, Cambridge Univ. Press,
o.A. Ivanova 1961.
D.M. Smirnov
AMS 1980 Subject Classification: 51 N1 0, 51 A03
AMS 1980 Subject Classification: 08A02
ASTRONOMETRY, MATHEMATICAL PROBLEMS
ASSOCIATOR of three elements a, b, c of a ring R - OF - Problems in astronomy related to the creation of
The element equal to an inertial reference coordinate system in space and a
(ab)c-a(bc). unification of the totality of fundamental astronomic
It is denoted by (a, b, c). constants by determining the coordinates of celestial
K.A. Zhevlakov bodies and by studying the rotation of the Earth.
AMS 1980 Subject Classification: 16-01,17-01 Methods in astronometry are based on the theory and
the practice of geometrical measurements on the celes-
ASTROID - A plane algebraic curve of order six, tial sphere and the determination of distances of stars.
described by a point M on a circle of radius r rolling An important element in astronometric studies is the
on the inside of a circle of radius R =4r; a hypocycloid most accurate possible determination of the direction of
with module m =4. Its equation in orthogonal Carte- the straight line 'observer - star' at the moment of the
sian coordinates is observation. Since the study of the mutual location of
x2/3+ y 2/3 = R2/3; points is more convenient and more illustrative than
the study of mutual location of directions, an auxiliary
and a parametric representation is spherical surface (the so-called celestial sphere) is intro-
X = R cos3 .!... y = R sin 3 .!... duced in the treatment, and all objects are assumed to
4' 4'
be at an equal distance from the observer and located
y
on this sphere. Spherical trigonometry makes it possible
to use various coordinate systems on the celestial
sphere and to determine numerous relations between
the angles and the arcs of configurations of celestial
objects. These relations determine a large proportion of
astronometric methods of observation and the geometr-
ical foundation of astronometric telescopes.
There are four cusps (see Fig.). The length of the arc The phenomena studied in astronometry are related

262

to various irregularities in the rotation of the Earth (the
irregular rotation of the Earth, migration of poles, vari-
ation of the rotation axis due to precession and nuta-
tion, etc.) on one hand, and to the actual displacements
of celestial bodies on the other. These considerations
determine the temporal and the qualitative structure of
the series of observations which are performed and
necessitate the use of specific mathematical methods for
their analysis. Hence the importance and the urgency
of solving and studying large systems of linear equa-
tions involving thousands of equations with hundreds
of unknowns, which is complicated even further by the
fact that some of them can only be insufficiently deter-
mined. A typical example of the genesis of such prob-
lems is the processing of observations made on the
bodies of the solar system (mainly small planets) in
order to determine the origin of the fundamental coor-
dinate system. Another example is the derivation of the
constituents of the gravitation field of the Earth from
the observations of the Earth made by artificial satel-
lites, realized by the worldwide network of observation
stations.
The complexity of the rotation of the Earth, the
difference between the real Earth, which is a visco-
elastic body, from the model as an absolutely solid
makes it necessary to search for latent relationships by
methods such as correlation-spectral analysis and vari-
ous smoothing methods. Such studies include the
studies of the variation in latitudes and longitudes and
the resulting pole migration, as well as evaluations of
fine and complicated fluctuations in the rate of rotation
of the Earth.
Owing to the nature of the observed phenomena it is
often impossible to complete the observations in a sin-
gle cycle, and it is then necessary to work with partly
overlapping segments of observations and with a vari-
able origin. Such observations are processed by
methods of solving finite-difference equations. This is
typical of instrumental research in astronometry and of
the problems in astronometry which involve the meas-
urement of some linear combinations of errors in the
coordinates of celestial objects.
Methods of mathematical statistics (mostly linear
methods) are employed to allow for errors in astro-
nometric observations.
v. V. Nesterov
V. V. Podobed
AMS 1980 Subject Classification: 85-01
ASTRONOMY, MATHEMATICAL PROBLEMS OF -
Mathematical problems in the study of celestial objects.
Special methods have been developed for the solution
of a number of such problems. These methods also
found use in other branches of science. On the other
ASTRO!'JOMY. MATHEMATICAL PROBLEMS OF
hand, mathematical techniques specially intended for
Earth-bound problems, which may be modified as
required, are extensively used in astronomy.
Astronomy is a complex science which deals with
celestial objects and their systems from various aspects,
which may be very remote from one another. This is
why the number of mathematical problems in astron-
omy is very large.
A very important chapter of astronomy is
astrometry, one of the main problems in which is to
determine an inertial reference coordinate system in
space. The coordinate systems which are traditionally
employed in astronomy, geodesy and other fields of sci-
ence, based on the planarity of the Earth's equator and
the vernal equinox (i.e. the intersection line of the
Earth's equator with the plane of the ecliptic), are not
inertial and cannot be firmly fixed in space, since both
these planes continuously execute a complex motion (as
a result of precession, nutation and migration of the
Earth's poles). For the sake of comparability, the coor-
dinates of the stars are usually referred to by the posi-
tion of the plane of the equator and of the point of ver-
nal equinox at some fixed date (epoch), while the coor-
dinate system itself, which has been fixed in this way, is
based on very careful determinations of the coordinates
of a number of stars, recorded in special catalogues
(fundamental star catalogues). However, a substantial
difficulty remains: In order to reproduce such a coordi-
nate system at a moment of time other than the epoch
of the catalogue, one must know the changes in the
locations of the fundamental stars relative to the coor-
dinate system, due to their proper motion. In order to
overcome this difficulty, it has been attempted since the
middle of the 20-th century to determine the inertial
coordinate system with respect to remote galaxies, the
proper motion of which is negligible. In this context
special importance is given in astronometry to prob-
lems involved in the calculation of the most probable
values of the parameters which determine the directions
of the stars on the strength of repeated observations,
and to estimating the probability characteristics of
these values. The solution of this problem is also
important in most other branches of astronomy, since
the latter is to a large extent an observational science.
Various mathematical problems must be solved in
theoretical astrophysics, which on the basis of observa-
tions of celestial objects studies the structure of these
objects, the physical processes taking place in them and
their evolution. One of principal problems of astrophy-
sics is the structure and evolution of stars. The theory
of the internal structure of stars leads to differential
equations which describe the conditions of mechanical
and energetic equilibrium of the stars. The solutions of
these equations can sometimes be expressed in terms of
263
ASTRONOMY, MATHEMATICAL PROBLEMS OF
elementary functions, but most often they are so com-
plicated that they must be solved by numerical
methods.
The study of atmospheres of stars and of processes
taking place in the nebulae and in the interstellar space
is based on the mathematical theory of radiative
transfer, which was intensively developed in astrophy-
sics. In certain cases, e.g. in the study of the passage of
radiation through a plane material layer, the transfer
equation reduces to integral equations, the solution of
which permits one to determine the characteristics of
the radiation field inside the star, and also of the radia-
tion emitted by the (interstellar) medium and accessible
to measurements.
In studying the motion of gas masses in stars and
nebulae, in studying the processes involved in the
extension of gas clouds, their collisions with each other
and with the interstellar space, extensive use is made of
the mathematical apparatus of gas dynamics and elec-
trodynamics.
Stellar astronomy, which is concerned with the laws
governing the structure, dynamics and evolution of stel-
lar systems, makes use of mathematical relationships
between the distribution of some given characteristics
of the stellar system (so-called distribution functions)
and the distribution of the observed characteristics. For
instance, the study of the connections (under certain
supplementary assumptions) between the distribution
functions of the distances of stars in a given solid angle
and their absolute and apparent (observational) values
yields an integral equation, the solution of which offers
an explanation of the law governing the distribution of
the stellar density in this solid angle. Similar equations
result on comparing the distribution functions of the
sought-for spatial velocities of the stars and the
observed radial velocities.
In stellar kinematics the problem of determining the
components of the Sun's velocity and the rotational
characteristics of the Galaxy from statistical studies on
the coordinates, the proper motion and the radial velo-
city, leads to an overdetermined system of constrained
equations compiled for individual stars (or individual
areas in the sky). The mathematical techniques in
mechanics are used in solving problems of stellar
dynamics connected with the study of star clusters,
galaxies and clusters of galaxies. The individual objects
which constitute the system are regarded as material
points which interact in accordance with the law of
gravitational attraction, with allowance for the specific
features of celestial bodies. A solution of the problems
of celestial mechanics, which is concerned with the
motion of celestial bodies in a gravitation field, leads to
systems of differential equations of motion. The solu-
tion of the most general problem of the motion of n
264
bodies which attract each other according to the law of
gravity under arbitrary initial conditions is solved by
numerical integration. However, the solutions given by
this method are only satisfactory during a limited
period of time, and no conclusions as to the evolution
of the system of objects are possible. The partial prob-
lem of three bodies was more satisfactorily solved using
power series in the powers of time; however, these
series converge very slowly and are therefore unsuitable
for application to the observations on the motion of
actual objects. Particular problems - the 3-body prob-
lem, the lO-body problem (the Sun and the 9 major
planets), etc. - have also been studied.
Problems involving the motion of specific celestial
bodies are solved by series expansions with respect to
powers of a small parameter and subject to various
assumptions that simplify the calculation.
Astrodynamics, which deals with the motion of artifi-
cial celestial objects, makes use of specific differential
equations of motion. Solutions of problems involved in
the motion of artificial satellites must allow for the dis-
turbing forces due to the non-spherical shape of the
Earth, the resistance of the atmosphere, the solar radia-
tion pressure (in the case of balloon satellites) and for
certain other factors.
For more details see the following entries:
Astrometry, mathematical problems of; Astrophysics,
mathematical problems of; Stellar astronomy, mathemat-
ical problems of; aassical celestial mechanics,
mathematical problems in.
N.P. Erpylev
AMS 1980 Subject Classification: 85-01
ASTROPHYSICS, MATIlEMATICAL PROBLEMS
OF - A type of problems in theoretical astrophysics in
which mathematical methods of investigation are exten-
sively employed. The principal task of theoretical astro-
physics is the interpretation of results of observations
with the purpose of studying the structure of the
objects encountered in the Universe, and to study the
physical processes taking place in it. In addition, there
are certain mathematical problems of a general nature
which are important not only in astrophysics, but also
in other branches of physics and in mathematics. The
most obvious example of the influence of mathematical
problems in astrophysics on the development of
mathematics is the so-called principle of invariance (cf.
Invariance, principle of) [1], which was first considered
in astrophysics, after which it was used to solve a wide
class of problems in mathematical physics and certain
problems in probability theory.
One of the principal problems in astrophysics is the
study of the structure and the evolution of stars, since
the bulk of the matter present in the Universe is con-

centrated in the stars. It is common to deal with the
internal structure of the stars separate from the struc-
ture of stellar atmospheres. The problems in the theory
of the internal structure are formulated as differential
equations, resulting from the conditions of mechanical
and energetic equilibrium of the star [2]. The condition
of mechanical equilibrium of a star is expressed as the
equality of the force of attraction directed towards the
interior of the star and the force of the pressure exerted
by the gases and by light, which is directed outwards.
The condition of energetic equilibrium is an expression
of the equality of the amount of energy produced in a
given volume of the star and the amount of energy
emitted by this volume. In addition to the equilibrium
conditions, the equations of state of the stellar matter,
the absorption coefficient of the radiation and the law
of emission of energy as a result of thermonuclear reac-
tions, as a function of the density of the matter and of
the temperature for the given chemical composition, are
assumed to be known. Owing to the complexity of the
equations, they are usually solved by numerical
methods. The results of the theory of stellar atmo-
spheres are sometimes employed as the boundary con-
ditions, which renders the conclusions more reliable.
These equations of stellar equilibrium may be reduced
in a significant number of cases to the Emden equation:
_1 -.!L
x 2 dx
[X2~]
dx
+yn = 0 (1)
with the conditions y = 1 and y' = 0 when x = O. If n
assumes one of the values 0, 1,5, the solutions of equa-
tion (1) may be expressed by elementary functions.
Equation (1) is applicable, for example, in the zones of
the star in which the energy transfer mainly takes place
by a convective flow, while the transfer of radiation
energy is relatively small. Another difficulty in develop-
ing the theory of the internal structure of stars is that
the interiors of stars cannot be observed. Observational
data are usually compared only with the calculated
integral characteristics of the star such as its mass, its
radius and the integral amount of energy emitted by
the star in unit time. The use of statistical material con-
cerning these characteristics makes it possible to
present hypotheses on the possible evolution of stars.
The theory of stellar atmospheres is based on the stu-
dies of processes of radiative transfer (cf. Radiative
transfer theory). The energy generated in the interior of
the star undergoes a complex process of transfer
through the atmosphere, after which it emerges to the
outside. The transfer equation is solved in this case
under conditions of radiation equilibrium, which
expresses the fact that each volume element of the
atmosphere radiates and absorbs equal amounts of
energy. The presence of a local thermodynamic equili-
ASTROPHYSICS. MATHEMATICAL PROBLEMS OF
brium is also frequently postulated. In theory, the solu-
tion of the transfer equation under the above condi-
tions yields the spectrum of the star, the comparison of
which with the results of observations provides infor-
mation on the structure of the stellar atmosphere and
on the processes taking place in it. In particular, the
determination of the chemical composition of stellar
atmospheres is based on the study of the line spectra of
the star.
Radiative transfer processes also take place in nebu-
lae and in the interstellar space. The transfer process
often consists in repeated light scattering as it passes
through matter. Sunlight, after being dispersed in the
planetary atmosphere, provides information about its
structure and its physical properties. Anisotropic light
scattering takes place in atmospheres of planets and in
dusty nebulae.
The theory of radiative transfer is one of the most
important branches of theoretical astrophysics. This
theory owes its impressive development [3], [4] to the
solution of the astrophysical problems mentioned
above. Its methods are also extensively employed in
other branches of physics such as geophysics (cf. Ge0-
physics, mathematical problems in), in the theory of
neutron transfer and in the calculations of the lumines-
cence of plasmas. Most frequently, radiation takes
place through a plane layer of matter. An example of
this are stellar and planetary atmospheres, since their
thickness is usually small in comparison to their radius.
In this case, for isotropic scattering, the radiative
transfer equation is reduced to the integral equation
y(x} ="2AXl K( I x-x ' I}y(x}dx
" +f(x}, (2)
where f (x) defines the distribution of the radiation
sources in the layer, while y (x) is the unknown radia-
tion field. The kernel K (x) of the integral equation and
the parameter "11. include the law of interaction of radia-
tion with matter, while Xo is the thickness of the layer,
expressed in unit lengths of the free path of radiation.
The study of integral equations of the type (2) is one of
the principal mathematical problems of radiation
theory. The mathematical problems of astrophysics
include the search for exact analytic solutions of such
equations [4]. In particular, if Xo = 00, "11.= 1 and
f(x) 0, as well as K(x)=E 1(x), where E 1(x) is the
first integral exponential function, equation (2) is
known as Milne's equation. Its solution (cf. Milne prob-
lem) under the assumption of certain supplementary
conditions yields the distribution of the temperature in
the interior of the stars. In more complicated problems
(anisotropic scattering, scattering of polarized light,
light scattering in a medium of arbitrary shape, etc.)
the resulting unknown function is not the function
265
ASTROPHYSICS, MATHEMATICAL PROBLEMS OF
y(x), but functions which depend not only on the coor-
dinates but also on the direction and on other magni-
tudes which characterize the radiation field at the given
point. These functions are determined with the aid of
appropriate systems of integral equations, which are a
generalization of (2).
These methods are used to find the complete solution
of the problem, which makes it possible to find both
the radiation field inside the medium and the charac-
teristics of the radiation emitted by the medium. In the
applications of the theory it is often only necessary to
know the emitted radiation. It has been shown that it
is possible to determine the characteristics of the radia-
tion emitted by the medium directly, without previously
finding the complete solution. This often simplifies the
task considerably, and is realized either by using the
integral equation (2) or by applying the principle of
invariance. The other approach consists in representing
the given layer as a sum of two layers and establishing
the relations between the characteristics of the radiation
emitted by the medium as a whole and the radiation
characteristics at the boundaries of the two layers.
These techniques are also applied in the theory of neu-
tron transfer, where they are known as the Albedo
method.
Another prominent subject in astrophysics is the
application of methods of gas dynamics [5] and electro-
dynamics [6] to the study of stars and interstellar space.
These methods are indispensable in the study of the
motion of gas masses in stars and nebulae. The hydro-
dynamic effects (cf. Hydrodynamics, mathematical prob-
lems in) in the atmospheres of the Sun and other stars,
the expansion of gas clouds and their collisions with
each other and with the more rarefied interstellar
medium, as well as resulting shock waves and turbulent
motions in it are some of the problems studied.
One of the more important features of interstellar gas
dynamics is the necessity of considering the interaction
of an ionized gas with magnetic fields. This interaction
affects both the motion of the gas and the changes in
the configuration and the energy density of the mag-
netic field. When the gas is in motion, its particles
'adhere' all the time, as it were, to the magnetic line of
force, moving along it or dragging it with them as they
move across the field. One says that that the lines of
force are 'frozen' into the matter (cf. Frozen-in
integral). The equations of motion of interstellar mag-
netic gas dynamics include: the continuity equation
(the law of conservation of mass); the induction equa-
tion of the magnetic field, which expresses the frozen-in
principle; the equation of the influx of the energy of
interstellar radiation (the law of conservation of
energy); and the equation which expresses the law of
conservation of momentum. The system is established
266
using the electrodynamic equations (Maxwell equations)
and the hydrodynamic equations. The solution of this
system is highly complicated, since the equations are
non-linear. There usually exist simplified variants, e.g.
one-dimensional and auto-modelling motions (cf.
Magneto-hydrodynamics, mathematical problems in).
Results obtained by radio-astronomic methods are
highly important in the development of astrophysics.
Since the radiation of radio waves usually takes place
in an ionized gas (a plasma), there arise the problems
of determining the properties of the plasma under the
conditions that exist in stars and in interstellar space
[7]. The behaviour of the plasma can be most fully
described by means of a kinetic equation. The differ-
ence between the kinetic equation for a plasma and the
kinetic Boltzmann equation for an ordinary gas consists
in the fact that the Coulomb interaction of the charged
particles of the plasma extends to relatively large dis-
tances, while in a gas consisting of neutral atoms and
molecules, the interactions are significant only during
direct collisions.
The study of the distribution of matter in the Galaxy
is made more difficult by the absorption of light by the
interstellar dust particles. This effect is particularly
marked in directions close to the Milky Way, since the
layer of interstellar dust is concentrated near the plane
of the Galaxy. This effect of light absorption is also
responsible for the decrease in the number of observed
galaxies as the direction of viewing approaches the
plane of our Galaxy. However, the absorbing substance
is not uniformly distributed; rather, it forms discrete
clouds of differing magnitudes, randomly distributed in
space. Thus there arises the problem of the statistical
study of the apparent distribution of the brightness of
the Milky Way and of the apparent distribution of
galaxies. These problems were studied in detail [1] in
the theory of fluctuations in brightness of the Milky
Way. The relevant equations were derived using the
principle of invariance. The mathematical problems of
the theory of the fluctuation of brightness are close to
the theory of certain types of random processes. In
their simplest variant it is assumed that the entire equa-
torial plane of the galaxy is infinitely filled, at a uni-
form density, with stars and with light-absorbing
clouds. If it is assumed that the proportion of the
transmitted light q is the same for all clouds, and that
g(u)du is the probability of the dimensionless bright-
ness being included between u and u + du, the function
g(u) will be given by the equation
g(u)+g ,(u) = -qg u }.
1 [ -q (3)
Equations of a similar type result on considering more
general models as well. The results of the observations

are frequently interpreted by using moments of the
magnitudes under study; accordingly, the solution of
the complicated initial equations is not invariably
necessary. The relationships of the moments, derived
from the initial equations, are quite simple and rapidly
yield the desired results.
The study of the structure, dynamics and evolution
of stellar systems forms the subject of statistical
mechanics of stellar systems. Here, systems of bodies
are usually considered as material points interacting in
accordance with the law of mutual attraction. The most
important systems are star clusters, galaxies and clus-
ters of galaxies. The given problems lead to questions
which are studied in physical kinetics. However, as dis-
tinct from studies of ordinary gases, the study of a stel-
lar gas involves very significant characteristic features.
First, a stellar gas is not at a complete statistical equili-
brium. Secondly, it is important to note that the gas is
situated in its own gravitation field, which should be
determined together with other unknown magnitudes.
Finally, the interaction between stars differs from the
interaction between gas molecules in that it concerns a
remote action discharge, which is analogous to the
Coulomb force acting between the charged particles in
a plasma. These features seriously complicate the
theoretical considerations. However, in studying the
motion of stars in a galaxy it is permissible to neglect
the interaction at a first approximation, since the relax-
ation times prove to be longer than the period of time
which determines the time of evolution of the Universe
[1], [8].
For star clusters the relaxation time is relatively
short, and it is then necessary to consider the interac-
tion between the stars. Here the usual practice is to
take advantage of the analogy between the Coulomb
forces and the gravitational forces and to use Landau's
kinetic equation. This equation was obtained for a
plasma using the Fokker- Planck approximation (cf.
Fokker- Planck equation) in representing the diffusion
process in the velocity space. Landau's equation cannot
give a sufficiently rigorous description of the process of
evolution of a stellar system, which inevitably involves
star losses. Accordingly, more general kinetic equations
are employed in studying the evolution of star clusters.
One of the more important stages of astrophysical
studies is to obtain the most complete and reliable
information possible about the objects being studied by
comparing the results obtained by observations with
theoretical results. The usual practice is to study a cal-
culated theoretical model and to check the validity of
the assumptions on which its construction is based. In
certain cases model representations can be dispensed
with and the observed characteristics can be directly
related to the theory, while reducing the number of jus-
ASTROPHYSICS, \IATHEMA TICAL PROBLEMS OF
tifiable initial assumptions to a minimum. A classical
example of such an approach is the problem of finding
the distribution of the spatial density of the stars from
their distribution in a projection [8]. The connection
between the observed and the spatial density distribu-
tion is obtained in the form of an Abel integral equa-
tion, the solution of which yields the desired result. The
same equation is also obtained in a number of other
problems, e.g. in the study of the distribution of atoms
with respect to altitude in the chromosphere and of the
electrons in the corona of the Sun. In all such cases
only one assumption - on spherical symmetry - is
made. A more complex integral equation is obtained
in the derivation of the distribution function of the spa-
tial velocities of the stars from the distribution of the
observed radial velocities. Here, the function is repro-
duced in three-dimensional space from the values of its
integrals in all possible planes. The examples here con-
sidered refer to so-called inverse problems. Inverse
problems in astrophysics involve considerable difficul-
ties and uncertainties both of mathematical nature and
in obtaining the necessary results of observations. The
mathematical aspects, in particular, have to do with
possible ill-posedness of the problems. The method of
regularization of the solutions of ill-posed problems
proves valuable in solving many astrophysical ques-
tions.
An example of an ill-posed inverse mathematical
problem in astrophysics is finding the optical properties
of the atmospheres of planets by measuring their pho-
tometric characteristics. Repeated scattering of light in
the atmosphere and reflections from the surface of the
planet result in a loss of information on the magnitudes
which determine the unknown optical properties of the
atmosphere. To obtain a single-valued solution, auxili-
ary simplifying assumptions must be made, which
strongly limit the possibilities of the investigation.
Inverse mathematical problems in astrophysics often
involve the solution of an integral equation
+00
F(x) = f A(x-y)j(y)dy (4)
of the first kind, which determines the transition from
the function F(x), which represents the known data
about some magnitude, obtained by observations, to
the unknown function f (x), which gives the true values
of this magnitude. The results of observations may be
distorted for various reasons - a fact which is quanti-
tatively described by the function A (x) in equation (4),
which is usually known for a given problem. For
instance, the functions F (x) and f (x) determine the
experimental and the true profile of the spectral line,
while the function A (x) corresponds to the observed
broadening profile of the ideally monochromatic spec-
267
ASTROPHYSICS, MATHEMATICAL PROBLEMS OF

tralline, caused by the properties of the spectrograph. without an orientation-reversing homeomorphism.

Many inverse problems arise in radio-astronomy Thus, for instance, the complex projective plane is an
both in processing the results of observations and in asymmetric variety, since the self-intersection of the
the interpretation of the data thus obtained. The pas- complex straight line is + 1 or - 1, depending on the
sage from the observed distribution of temperature to orientation. Certain knots can differ from their mirror
the true distribution in the one-dimensional case is image owing to the fact that their branched coverings
effected by the use of equation (4), in which A (x) are asymmetric varieties.
A. V. Chernavskif
determines the property of the radio-telescope antenna
and is called its diagram. An important feature of Editorial comments. This notion can be found, e.g., in
radio-astronomical studies is the use of the 21-cm radio [A1], Chapt. 5.
wave-length, which is generated by the transition References
[A1] HIRSCH, M.W.: Differential topology, Springer, 1976.
between the sublevels of the superfine structure of the
hydrogen atom. The study of the interstellar radio-type AMS 19aO Subject Classification: 57MXX
radiation on this wave-length yields valuable informa-
tion on the structure of the Galaxy and its rotation, on ASYMMETRY COEFFICIENT - The most frequently
the distribution of the interstellar gas in it, and on the employed measure of the asymmetry of a distribution,
structure of the central part of the Galaxy which is defined by the relationship
inaccessible for optical methods. The line profiles _ IL3
observed in the various directions have a complicated "Yl - 3'72'
IL2
shape due to various factors (thermal motion of atoms,
where JL2 and J.L3 are the second and third central
chaotic motion of the interstellar gas, the rotation of
moments of the distribution, respectively. For distribu-
the Galaxy, etc.). Accordingly, a number of inverse
tions that are symmetric with respect to the mathemati-
problems arise on passing from the measured to the
sought characteristics. In order to eliminate the blurring cal expectation, YI =0; depending on the sign of YI one
effect of thermal and chaotic macroscopic movements speaks of positive asymmetry (YI >0) and negative asym-
metry (YI <0). In the case of the binomial distribution
of the gas from the observed profile, equation (4) is
corresponding to n Bernoulli trials with probability of
employed, in which A (x) is a description of the velocity
success p,
distribution. This routine systematically yields gas velo- _ i-2p
(*)
Y1 - -Y"'n=rp=(}:'_=p=) ,
cities and other magnitudes.
References one has: If p = 1 /2 (YI =0), the distribution is sym-
[l) AMBARTSUMYAN, V.A.: Scientific work, Vol. I, Erevan, 1960 (in metric; if p < 1 / 2 or p> 1 / 2, one obtains typical dis-
Russian). tribution diagrams with a positive (Fig. a) and negative
[2) ScHWARZSCHILD, M.: Structure and evolution of stars, Prince-
ton Univ. Press, 1958. (Fig. b) asymmetry.
[3) SOBOLEV, V.V.: A treatise on radiative transfer, v. Nostrand,
1960 (translated from the Russian).
[4) SOBOLEV, V.V.: Course in theoretical astrophysics, NASA,
Washington, D.C., 1969 (translated from the Russian).
[5) KAPLAN, S.A.: Interstellar gasdynamics, Moscow, 1958 (in Rus-
sian).
[6) PIKEL'NER, S.B.: Fundamentals of cosmic electrodynamics, Mos-
cow, 1966 (in Russian).
a) P(k, 10, 1/5)
I.N. Minin
Editorial comments.
References
[A1] EHLERS, J. (ED.): Relativity theory and astrophysics, 1-3,
Amer. Math. Soc., 1967.
[A2) KURT, R.: Introduction to stellar statistics, Pergamon Press, k
1967.
[A3] LANDSBERG, P.T. and EVANo, D.A.: Mathematical cosmology,
Clarendon Press, 1977. b) P(k, 10,4/5)
[A4) CHANDRASEKHAR, S.: Radiative transfer, Dover, reprint, 1960.
[AS] CHANDRASEKHAR, S.: Stellar structure, Dover, reprint, 1957.
[A6) HAWKING, S.W. and ELLIS, G.F.R.: The large scale structure Diagrams of the binomial distribution
of space-time, Cambridge Univ. Press, 1973. P (k, n, p) corresponding to n = 10 Ber-
AMS 1980 Subject Classification: 85-01 noulli trials: a) with positive asymmetry
(p = 1/5); b) with negative asymmetry
ASYMMETRIC VARIETY - An oriented variety M (p=4/5).

The asymmetry coefficient (*) tends to zero as n~oo,
in accordance with the fact that a normalized binomial
distribution converges to the standard normal distribu-
tion.
The asymmetry coefficient and the excess coefficient
are the most extensively used characteristics of the
accuracy with which the distribution function Fn(x) of
the normalized sum
(X, + ... + Xn)- nIL,
y;;;;
where XI, ... ,Xn are identically distributed and mutu-
ally independent with asymmetry coefficient 8 1, may be
approximated by the normal distribution function
CP(x) = . b-
v2'IT
Je-,2 12 dz.
-x
Under fairly general conditions the Edgeworth series
yields
Fn(x) = CP(x)- Fn ~ cp(3)(x)+O l~l'
where qP)(x) is the derivative of order three.
References
[I] CRAMER, H.: Mathematical methods o/statistics, Princeton
Univ. Press, 1946.
[2] WILKS, S.S.: Mathematical statistics, Wiley, 1962.
A. V Prokhorov
Editorial comments, The asymmetry coefficient is usu-
ally calied the coefficient of skewness. One correspondingly
speaks of the skewness of a distribution and of positive,
respectively negative, skewness.
The excess coefficient is more often called the coefficient
of kurtosis.
AMS 1980 Subject Classification: 62-XX
ASYMMETRY OF A DISTRIBUTION - A qualita-
tive property of the curve of a distribution, indicating
that the distribution deviates from a symmetric one.
The asymmetry of a distribution is positive (negative) if
the asymmetry coefficient is positive (negative). If the
asymmetry of a distribution is positive (negative), then
the 'longer' part of the density curve of the distribution
is found to the right (left) of the mode.
K.P. Latyshev
Editorial comments. In Western literature one commonly
calis this notion the skewness of a distribution, and accord-
ingly one speaks of the skewness coefficient.
AMS 1980 Subject Classification: 60EXX, 62EXX
ASYMPTOTE of a curve y = f (x) with an infinite
branch - A straight line the distance of which from the
point (x, f (x on the curve tends to zero as the point
moves along the branch of the curve to infinity. An
ASY.\1PTOTIC BASIS
asymptote can be vertical or inclined. The equation of
a vertical asymptote is x = a, where f (x)~ + 00 ( - 00)
as x~a (from one side) is satisfied. An inclined
asymptote, with equation y = kx + I, exists if and only if
the limits
k = limM., I = lim[f(x)-kx),
x
exist as x~+oo (or as x~-oo).
Similar formulas are also obtained for parametrized
(unbounded) curves in general parametric representa-
tion. In polar coordinates an asymptote of a curve
r = r( cp), where r >0, with slope angle lX, is defined by
the condition r~ + 00 as cp~a. The distance p of this
asymptote from the coordinate origin is calculated by
the formula
p = lim It Ir(a+t) as t~+O (or as t~-O).
If there exists a limit position of the tangent line to
the infinite branch of the curve, this position is an
asymptote. The converse is not always true. Thus, the
curve y = (sin x 2) / x has the asymptote y = as
x~-+- 00, even though a limit position of the tangent
line does not exist. Hyperbolas are the only second-
order curves with asymptotes. The asymptotes of the
hyperbola (x 2 / a 2 )_(y2 / b 2 )= 1 are given by the
equations (x / a)-+-(y / b)=O. An inclined asymptote
yields a simple (linear with respect to x) asymptotic
approximation of the function:
f(x) = kx+l+o(l)
as x~+oo (or as x~-oo).
References
[I] RASHEVSKIi, P.K.: A course 0/ differential geometry, Moscow,
1956 (in Russian).
[2] KUDRYAVTSEV, L.D.: Mathematical analysis, Moscow, 1973 (in
Russian).
L.P. Kuptsov
Editorial comments.
References
[A 1) POGORELOV, A.V.: Differential geometry, Noordhoff, 1959
(translated from the Russian).
AMS 1980 Subject Classification: 53A04, 26A 12
ASYMPTOTIC BASIS, asymptotic basis of order k -
A sequence of natural numbers and zero, which as a
result of its summation repeated k times yields all suffi-
ciently large natural numbers. The number k is called
the order of the asymptotic basis. Thus, the sequence of
prime numbers is an asymptotic basis of order 4 (LM.
Vinogradov, 1937); the sequence of cubes of natural
numbers is an asymptotic basis of order 7 (Yu.V. Lin-
nik, 1942).
B.M. Bredikhin
Editorial comments. See also Waring problem; Gold-
bach problem.
AMS 1980 Subject Classification: 10105
269
ASYMPTOTIC DENSITY
ASYMPTOTIC DENSITY - A variant of the general
concept of the density of a sequence of natural
numbers, which measures how large a part of the
sequence of all natural numbers belongs to the given
sequence A of natural numbers including zero. The
asymptotic density of a sequence A is expressed by the
real number a defined by the formula
I . A (x\
a = lmmf~,
x------.oc X some neighbourhood of the point Xo (except possibly at
where
A(x) = ~ I, x;;;;.l.
UEA.
O<a~x
The number
f3 = lim sup~
x I.e.
x----"'oo
is known as the upper asymptotic density. If the numbers
a and f3 coincide, their common value is called the
natural density. Thus, for instance, the sequence of
numbers that are free from squares has the natural den-
sity () = 6/"? The concept of an asymptotic density is
employed in finding criteria for some sequence to be an
asymptotic basis.
B. M. Bredikhin
Editorial comments. The number a as defined above is
also called the lower asymptotic density.
References
[A 1] HALBERSTAM, H. and ROTH, K.F.: Sequences, 1, Clarendon
Press, 1966.
AMS 1980 Subject Classification: 10L02
ASYMPTOTIC DERIVATIVE - A
approximate derivative.
AMS 1980 Subject Classification: 26A24
ASYMPTOTIC DIRECTION - A direction on a regu-
lar surface in which the curvature of the normal section
of the surface is zero. For the direction du: dv at a
point P to be an asymptotic direction, the following
condition is necessary and sufficient:
Ldu 2 +2Mdudv+Ndv 2 = 0,
where u and v are interior coordinates on the surface,
while L, M and N are the coefficients of the second
fundamental form of the surface, calculated at P. At an
elliptic point of the surface the asymptotic directions
are imaginary, at a hyperbolic point there are two real
asymptotic directions, at a parabolic point there is one
real asymptotic direction, and at a flat point any direc-
tion is asymptotic. Asymptotic directions are self-
conjugate directions (d. Conjugate directions).
References
[II RASHEVSKli, P.K.: A course of differential geometry, Moscow,
1956 (in Russian).
Editorial comments.
270
References
[A1] HsIUNG, c.c.: A first course in differential geometry, Wiley,
1981.
[A2] STRUIK, 0.1.: Lectures on classical differential geometry,
Addison-Wesley, 1950.
AMS 1980 Subject Classification: 53A05
ASYMPTOTIC EQUALITY - Two functions f (x)
and g(x) are called asymptotically equal as X-Ho if in
x itself)
j(x) = (x)g(x),
where
=
lim (x) I,
X-Xo
f(x) = g(x)[1 +0(1)),
as X-Ho (xo is a finite or an infinite point of the set
on which the functions under consideration are
defined). If g(x) does not vanish in some neighbour-
hood of xo, this condition is equivalent to the require-
ment
lim M.. = 1.
X~Xo g(x)
In other words, asymptotic equality of two functions
f(x) and g(x) as X-Ho means, in this case, that the
relative error of the approximate equality of f (x) and
g(x), i.e. the magnitude [f(x)-g(x)]/ g(x), g(x)=,FO,
is infinitely small as X-Ho. Asymptotic equality of
functions is meaningful for infinitely-small and
infinitely-large functions. Asymptotic equality of two
synonym for functions f(x) and g(x) is denoted by f(x)~g(x) as
X~Xo, and is reflexive, symmetric and transitive.
Accordingly, the set of infinitely-small (infinitely-large)
functions as X~Xo is decomposed into equivalence
classes of such functions. An example of
asymptotically-equal functions (which are also called
equivalent functions) as X~Xo are the functions u(x),
sin u(x), In(l +u(x)], eU(x) -I, where limHXo u(x)=O.
If f~fl and g~gl as X~Xo, then
lim 1.l!:l = lim flex) ,
,~\u g(x),~xu gl(X)
where the existence of anyone of the limits follows
from the existence of the other one. See also Asymptotic
expansion of a function; Asymptotic formula.
M.1. Shabunin
Editorial comments. One also says that f (x) and g(x)
are of the same order of magmtude at x(] instead of asymp-
totically equal.
References
[A 1] COURANT, R. Differential and integral calculus, 1, Blackie
and Son, 1948, Chapt. 3, Sect. 9.
AMS 1980 Subject Classification 26E99
E. V Shikin
ASYMPTOTIC EXPANSION of a function j(x) - A
 ASYMPTOTIC LINE

senes often disregarded if asymptotic approximations can be

relatively easily found.
n =0
References
such that for any integer N;;.O one has [I] POINCARE, H.: 'Sur les integrales irregulieres des equations
N lineaires', Acta Math. 8 (1886),295-344.
f(x) = 2:,"'n(X)+O(<PN(X (x-Ho), (I) [2] WHITTAKER, E.T. and WATSON, G.N.: A course of modern
n =0
analysis, Cambridge Univ. Press, 1952.
where {<Pn(x)} is some given asymptotic sequence as [3] ERDELYI, A. and WYMAN, M.: 'The asymptotic evaluation of
certain integrals', Arch. Rational Mech. Anal. 14 (1963),217-
X~Xo. In such a case one also has
oc
260. M. V. Fedoryuk
f(x) ~ 2:, "'n(x), {<Pn(x)}, (x-no) (2) Editorial comments.
n =0
References
The sequence {<Pn(x)} is omitted from formula (2) if it [A1] BRUUN, N.G. DE: Asymptotic methods in analysis, Dover,
is clear from the context which sequence is meant. reprint, 1981.
The asymptotic expansion (2) is called an asymptotic [A2] ERDELYI, A.: Asymptotic expansions, Dover, reprint, 1956.
expansion in the sense of Erdelyi [3]. An expansion of [A3] COPSON, E.T.: Asymptotic expansions, Cambridge Univ.
Press, 1965.
the type 00
[A4] MURRAY, J.P.: Asymptotic analysis, Springer, 1984.
f(x) ~ 2:, an<Pn(X) (x-Ho), (3) [A5] BLEISTEIN, N. and HANDELSMAN, R.A.: Asymptotic expan-
n =0
sions of integrals, Dover, reprint, 1986.
where an are constants, is called an asymptotic expan- AMS 1980 Subject Classification: 41 A60
sion in the sense of Poincare. If the asymptotic sequence
of functions {<Pn(x)} is given, the asymptotic expansion ASYMPTOTIC EXPRESSION - A synonym for
(3), contrary to the expansion (2), is uniquely defined asymptotic formula.
by the function f (x) itself. If (1) is valid for a finite AMS 1980 Subject Classification: 41 A60
number of values N=O, ... ,No<oo, then (1) is called
an asymptotic expansion up to O(<PNo(X. The series ASYMPTOTIC FORMULA - A formula containing
x 00 the symbols 0, 0 or the equivalence sign "-' (asymptotic
2:, "'Ax), 2:, an <Pn (X) equality of functions).
n =0 n =0
Examples of asymptotic formulas:
are known as asymptotic series. As a rule such series are sinx = x +o(x2), x~O;
divergent. Asymptotic power series are the ones most
commonly employed; the corresponding asymptotic COSX = 1+O(X2), x~O;

expansions are asymptotic expansions in the sense of x 3+x+1 ~ x3, x~oo;

Poincare. X
17(X) ~ Inx' x~oo
The following is an example of an asymptotic expan-
sion in the sense of ErdeIyi: ('1T(x) is the amount of prime numbers not exceeding x).
B.M. Bredikhin
I,(x) ~ -..Vf2[COS
-:;;; [x- 17V2 _!!....J
4
~(-I)na2nX-2n-
n =0
Editorial comments. On the meaning of the symbols
0,0 and ~, see e.g. [A1] or [A2].

-sin [x- 17{ - : t~0(-I)"a2n+lx-2n-l] References

(x~ + 00), where J vex) is the Bessel function, and
a = rev +n + I /2)
n 2n n!r(v-n+I/2)
The concepts of an asymptotic expansion of a func-
tion and of an asymptotic series were introduced by H.
Poincare [1] in the context of problems in celestial
mechanics. Special cases of asymptotic expansions were
discovered and utilized as early as the 18-th century [2].
Asymptotic expansions play an important role in many
problems in mathematics, mechanics and physics. This
is because many problems do not admit exact solutions,
but their solutions can be obtained as asymptotic
approximations. Moreover, numerical methods are
[A1] BRUUN, N.G. DE: Asymptotic methods in analysis, Dover,
reprint, 1981.
[A2] ERDELYI, A.: Asymptotic expansions, Dover, reprint, 1956.
AMS 1980 Subject Classification: 41 A60
ASYMPTOTIC LIMIT - A synonym for approximate
limit.
AMS 1980 Subject Classification: 26A24
ASYMPTOTIC LINE - A curve r on a regular sur-
face F such that the normal curvature along r is zero.
An asymptotic line is given by the differential equa-
tion:
II 2= Ldu +2Mdudv+Ndv
2 = O.
where II is the second fundamental form of the surface.

271
ASYMPTOTIC LINE
The osculating plane of an asymptotic line r, if it
exists, coincides with the tangent plane to F (at the
points of f), and the square of the torsion of an
asymptotic line is equal to the modulus of the Gaussian
curvature K of the surface F (the Beltrami- Enneper
theorem). A straight line I EF (e.g. a generating line of a
ruled surface) is always an asymptotic line. If r is a
parabolic curve (e.g. a circle on a standard torus
separating the domains of Gaussian curvatures of dif-
ferent signs), it is an asymptotic line.
A unique asymptotic line, which coincides with the
rectilinear generator, passes through each point of a
parabolic domain (where K=O, but n*O). Through
each point of a hyperbolic domain (where K <0) there
pass exactly two asymptotic lines, forming the so-called
asymptotic net, which plays an important role in the
study of the spatial form of a surface of negative curva-
ture (cf. Negative curvature, surface of). For instance,
on a complete surface this net is homeomorphic to the
Cartesian net on the plane if
grad I vi ~ K I ~ q, q = const.
Asymptotic nets on surfaces of constant negative curva-
ture are Chebyshev nets (cf. Chebyshev net), and the
surface area of a quadrangle formed by asymptotic
lines is proportional to the excess of the sum of its inte-
rior angles (Xi over 2w:
I K IS = 2'1T-0:] -0:2 -0:3 -0:4
(Hazzidakis' formula).
Under a projective transformation w of the space, the
asymptotic lines of a surface F become the asymptotic
lines of the transformed surface w(F).
Asymptotic lines on surfaces in a three-dimensional
Riemannian space are defined in a similar manner.
Various generalizations of the concept of an asymptotic
line on manifolds imbedded in a multi-dimensional
space are known; the most frequently-used one involves
the concept of the second fundamental form, which is
associated with a given normal vector.
References
[I] POGORELOV. A.V.: Differential geometry. Noordhoff. 1959
(translated from the Russian).
[2] RASHEVSKli, P.K.: A course of differential geometry, Moscow,
1956 (in Russian).
M.I. Voitsekhovski(
Editorial comments. Hazzidakis' formula can be found in
[A1] and [A2J, p. 204.
References
[A 1] HAZZIDAKIS. 1.N. 'Uber elnlge Eigenschaften der Flachen
mit konstanten Krummungsmasz', Crelle's J Math. 88
(1880), 68-73.
[A2] STRUIK. D.l.: Lectures on classical differential geometry,
Addison-Wesley, 1950.
[A3] HSIUNG. c.c.: A first course In differential geometry, Wiley,
1981.
272
[A4] SPIVAK, M.: Differential geometry, 3, Publish or Perish, 1975.
[AS] HICKS, N.J.: Notes on differential geometry, v. Nostrand,
1965.
[A6] BLASCHKE, W. and LEICHTWEISS, K.: Elementare Differential-
geometrie, Springer, 1973.
AMS 1980 Subject Classification: 53A05
ASYMPTOTIC NEGLIGIBILITY - A property of
random variables indicating that their individual contri-
bution as components of a sum is small. This concept
is important, for example, in the so-called triangular
array. Let the random variables Xnk (n = 1, 2, ... ;
k = 1, ... ,kn ) be mutually independent for each n, and
let
If for all i>O and S>O, at sufficiently large values of n,
the inequality
(1)
is .satisfied, the individual terms Xnk are called asymp-
totically negligible (the variables Xnk then form a so-
called zero triangular array). If condition (1) is met,
one obtains the following important result: The class of
limit distributions for Sn - An (An are certain 'center-
ing' constants) coincides with the class of infinitely-
divisible distributions (cf. Infinitely-divisible distribu-
tion). If the distributions of Sn converge to a limit dis-
tribution, kn----,>oo, and the terms are identically distri-
buted, condition (1) is automatically met. If the
requirement for asymptotic negligibility is strengthened
by assuming that for all {>o and S>O for all suffi-
ciently large n one has
P [max
l~k~k"
I Xnk I >] < 0, (2)
then the following statement is valid: If (2) is met, the
limit distribution for Sn - An can only be a normal dis-
tribution (in particular with variance equal to zero, i.e.
a degenerate distribution).
A. V. Prokhorov
Editorial comments.
References
[A 1] FELLER, W.: An introduction to probability theory and its
applications, 2, Wiley, 1966.
AMS 1980 Subject Classification: 60-XX
ASYMPTOTIC NET - A net on a surface formed by
two families of asymptotic lines (cf. Asymptotic line).
An asymptotic net only exists on non-developable sur-
faces of non-positive curvature. The orthogonality of an
asymptotic net characterizes a minimal surface.
V T Bazr/ev
AMS 1980 Subject Classification: 53A05
ASYMPTOTIC POINT - A singular point of a curve
around which it twists an infinite number of times,
 ASYMPTOTIC VALCE

approaching it without limit. For example, the asymp- uniformly in argz as I z I~oo inside any closed angle
totic point of the logarithmic spiral. contained in D, then
AMS 1980 Subject Classification: 51 N1 0 , ~ nan
f(z) ~ - ...-
n=1 zn+1
ASYMPI'OTIC POWER SERIES - An asymptotic
series with respect to the sequence
uniformly in arg z as I z I~ 00 in any closed angle con-
tained in D.
{x-n} (x~oo)
References
or with respect to a sequence [1) COPSON, E.T.: Asymptotic expansions, Cambridge Univ. Press,
1965.
{(x - xo)"} (x~xo) [2) ERDELYI, A.: Asymptotic expansions, Dover, reprint, 1956.
[3) WHITTAKER, E.T. and WATSON, G.N.: A course of modem
(cf. Asymptotic expansion of a function). Asymptotic analysis, Cambridge Univ. Press, 1952.
power series may be added, multiplied, divided and M.l. Shabunin
integrated just like convergent power series. Editorial comments.
Let two functions f(x) and g(x) have the following References
asymptotic expansions as x~oo: [A1] BRUUN, N.G. DE: Asymptotic methods in analysis, Dover,
x a x b reprint, 1981.
f(x) ~ ~~, g(x) ~ ~~.
n=OX n=OX AMS 1980 Subject Classification: 41 A60
Then
~ Aan+Bbn ASYMPTOTIC REPRESENTATION - A synonym for
I) Af(x) +Bg(x) ~ ... n
n=O X asymptotic formula.
(A, B are constants); AMS 1980 Subject Classification: 41 A60
ex; C
2)f(x)g(x) ~ ~~;
n=O X ASYMPTOTIC SEQUENCE - A sequence of func-
1 1 d tions {<Pn(x)} such that
3) - - ~ -+ ~_n ao=FO
00

f(x) ao n=1 xn ' <Pn+I(X) = o(<Pn(x, x~Xo, xEM,

(cn , dn are calculated as for convergent power series); where Xo is a limit point of the set M (finite or infin-
4) if the function f (x) is continuous for x >a >0, ite). If the nature of M is clear from the context, then
then ex; one simply writes X~Xo. If {<Pn(x)} is an asymptotic
jff(t)-ao-!!J...]dt ~ ~ an+nl ; sequence and !f,{x) is a function defined on M, then
xl t n=lnx N(x )<Pn(x)} will also be an asymptotic sequence.
5) an asymptotic power series cannot always be dif- Examples of asymptotic sequences:
ferentiated, but if f (x) has a continuous derivative 1) {(x-xo)n}, X~Xo;
which can be expanded into an asymptotic power 2) {x- n }, x~oo;
series, then 3) {eXx- n }, x~oo;
j'(x) ~ _ ~ (n-I~an-I. 4) {z-n}, z~oo, ZED, where D is an unbounded
n=2 X domain in the complex plane. Asymptotic sequences
Examples of asymptotic power series. such as 1), 2) and 4) are called asymptotic power
sequences.
OOje x - t ex; (_l)n-I(n-I)! M.l. Shabunin
-dt ~ ~ - - - -'
x t n= I xn ' Editorial comments.
Yxe-iXmp(x) ~ ~ e- l7T
/
4(-i)"[f(n+I/2)j2 References
n=O 2n-I/2~/2n!xn ' [A1] BRUUN, N.G. DE: Asymptotic methods in analysis, Dover,
reprint, 1981.
where Hb1) (x) is the Hankel function of order zero (cf. [A2] ERDEL YI, A.: Asymptotic expansions, Dover, reprint, 1956.
Hankel functions) (the above asymptotic power series
AMS 1980 Subject Classification: 41 A60
diverge for all x).
Similar assertions are also valid for functions of a
ASYMPI'OTIC SERIES - See Asymptotic expansion
complex variable Z as z~oo in a neighbourhood of the
of a function.
point at infinity or inside an angle. For a complex vari-
able 5) takes the following form: If the function f (z) is AMS 1980 Subject Classification: 41 A60
regular in the domain D = { I z I >a, a< I argz I <f3}
and if ASYMPTOTIC VALUE - A limit value along some
00 a
f(z) ~ ~~ path. More exactly, a complex number a or a= 00 is
n=O Z called an asymptotic value for a function f (z) of the

273
ASYMPTOTIC VALUE

complex variable z at a point a of the closure 15 of its observation. An estimator Yn E~ which attains the
domain of definition D if there exists a path L: lower bound just mentioned is asymptotically efficient.
z =z(t), O.;;;;t<l, LCD, terminating at a, i.e. so that Under certain conditions this property is satisfied by
lim z(t) ::::: a, the maximum-likelihood estimator for (), which makes
t~I-O
the classical definition meaningful. If the
along which asymptotically-efficient estimator Yn exists, the magni-
limf(z) ::::: a, ZEL.
z~a

tude rJ2 (Vn r..

For instance, at the point a = 00 the function f (z) = eZ lim n)
n~oo 02(Vn Tn)
has the asymptotic values a] = 0 and a2 = 00 along the
paths L]: z=-t, O.;;;;t<+oo, and L 2 : z=t, is called the asymptotic relative efficiency of Tn. Cer-
O';;;;t < + 00, respectively. Sets of asymptotic values play tain variants of the concept of an asymptotically-
an important role in the theory of limit sets (cf. Limit efficient estimator are due to R.A. Fisher, C.R. Rao
set). and others.
If f (z) has two different asymptotic values at a, a is References
called a point of indeterminacy for the function f (z). [1) RAo, C.R.: Linear statistical inference and its applications,
For any function fez), defined in a simply-connected Wiley, 1965.
O. V Shalaevskil
plane domain, the set of points of indeterminacy is at
Editorial comments. More modern definitions of this
most countable.
concept are due to J. Hajek, L. LeCam and others.
The above definition of asymptotic value refers to
References
asymptotic point values. If the limit set of a curve L is a [A 1] IBRAGIMov, J.A. and HAS'MINSKII, R.Z. [R.Z. KHAS'MINSKIT]:
set E C aD rather than a single point a EaD, one also Statistical estimation. Asymptotic theory, Springer, 1981
speaks of the asymptotic value A (j, E) associated with (translated from the Russian).
E. AMS 1980 Subject Classification: 62F12
References
[I] COLLINGWOOD, E.F. and LOHWATER, A.1.: The theory of clus- ASYMPTOTICALLY-STABLE SOLUTION - A solu-
ter sets, Cambridge Univ. Press, 1966, Chapt. I; 6.
tion of a differential system that is stable according to
[2] MAcLANE, O.R.: Asymptotic values of holomorphic functions,
Rice Univ. Studies, Math. Monographs, 49, Rice Univ., Hous- Lyapunov (d. Lyapunov stability) and that attracts all
ton, 1963. the other solutions with sufficiently close initial values.
VI. Gavrilov
Thus, the solution
E.D. Solomentsev
XC-T, ~o), x (lX, ~o) ::::: ~o,
Editorial comments. The most famous results on asymp-
totic values is the Oenjoy- Carleman- Ahlfors theorem. Let of the system dx
fez) be an entire function with n distinct (finite) asymptotic dT ::::: J(T, x) (*)
values at the point 00. Then fez) must be of order :?n / 2.
with a right-hand side f(T, ~), given for all T;;;'a, ~ERn,
This result was conjectured by A. Oenjoy (1907). The first
complete proof was given by L. Ahlfors (1929), after T.
and which is such that solutions of (*) exist and are
Carleman had obtained a less sharp result. See, for exam- unique, will be an asymptotically-stable solution if,
ple, [A1], Sect. 60. together with all its sufficiently close solutions
References X(T,~), I~-~o I <h, h>O,
[A 1] DING HAS, A.: Vorlesungen Uber Funktionentheorie, Springer,
1961.
it is defined for all T;;;'a and if for an arbitrary >0
there exists a 0, O<o<h, such that I ~-~o I <0 implies
AMS 1980 Subject Classification: 30D30, 30D40
II X(T, ~)- X(T, ~o) II < {
ASYMPTOTICALLY-EFFICIENT ESTIMATOR - A for all T;;;'a and
concept which extends the idea of an efficient estimator
to the case of large samples (cf. Effective estimator).
An asymptotically-efficient estimator has not been as T _____ + 00.
uniquely defined. Thus, in its classical variant it con- The concept of an asymptotically-stable solution was
cerns the asymptotic efficiency of an estimator in a introduced by A.M. Lyapunov [I]; it, together with
suitably restricted class R of estimators. In fact, let Tn various special types of uniform asymptotic stability, is
be a consistent estimator of a one-dimensional parame- extensively used in the theory of stability [2], [3], [4].
ter () constructed from a random sample of size n. Then References
Tn ESt if the variance rl-cVn Tn) exists,and if it is [1] LYAPliNOV, A.M.: 'ProbJeme general de la stabilite du mouve-
ment', in Ann. of Math. SlUdies, Vol. 17, Princeton Univ. Press,
bounded from below, as n-H/J, by the inverse of the 1947.
Fisher amount of information corresponding to one [2] KRASOVSKI1, N.N.: Stahility o(mOlion. Applications of

274

Lyapunov's second method to differential systems and equations
with delay, Stanford Univ. Press, 1963 (translated from the
Russian).
[3] HAHN, W.: Theorie und Anwendung der direkten Methode von
Ljapunov, Springer, 1959.
[4] RoveHE, N., HABETS, P. and LALOY, M.: Stability theory by
Liapunov's direct method, Springer, 1977. v S B d
1 U.. og anov
Editorial comments.
References
[A1) HAHN, W.: Stability of motion, Springer, 1967.
AMS 1980 Subject Classification: 34005, 34020
ASYMPTOTICALLY-UNBIASED ESTIMATOR - A
concept indicating that the estimator is unbiased in the
limit (cf. Unbiased estimator). Let XI, X 2 , . . . , be a
sequence of random variables on a probability space
(n, S, P), where P is one of the probability measures in
a family fJJ. Let a function g(P) be given on the fam-
ily fJJ, and let there be a sequence of S-measurable
functions Tn(X I , . . . ,Xn), n = 1, 2, ... , the mathemat-
ical expectations of which, Ep Tn (X I, . . . ,Xn), are
gJven. Then, if, as n .....H:fJ,
EpTn(Xj"",Xn)~g(P), PE9,
one says that Tn is a function which is asymptotically
unbiased for the function g. If one calls X I, X 2, . . . ,
'observations' and Tn 'estimators', one obtains the
definition of an asymptotically-unbiased estimator. In
the simplest case of unlimited repeated sampling from a
population, the distribution of which depends on a
one-dimensional parameter ()Ee, an asymptotically-
unbiased estimator Tn for g(), constructed with respect
to the sample size n, satisfies the condition
Eo Tn(X j , ... ,Xn) ~ g(8)
for any ()Ee, as n~oo.
0. V. Shalaevskil
AMS 1980 Subject Classification: 62F12
ASYMPTOTICALLY-UNBIASED TEST - A concept
indicating that the test is unbiased in the limit. For
example, in the case of n independent samples from a
one-dimensional distribution depending on a parameter
()En, let H be the null hypothesis: ()En H , and let K be
the alternative:
8 E S2 K, S2 H U S2K = Q, QH U QK = 0.
The critical set Rn in the n-dimensional Euclidean
space, n = 1, 2, ... , is an asymptotically-unbiased test
of the hypothesis H with level LX if
lim P(R n 18) .:;;; a, 8EQH,
n~x
a .:;;; lim P(R n 18), 8EQK'
n~x
The function
lim P(R n 18)
n~oo
ATOM
is called the asymptotic power function of the test Rn.
0. V. Shalaevskil
AMS 1980 Subject Classification: 62F05
ASYMPTOTICS OF ARITIIMETIC FUNCflONS,
asymptotics of number-theoretical functions - Approxi-
mate representations of arithmetic functions (functions
defined for all natural number values of the argument)
by means of comparatively simple expressions with
arbitrary small error terms. More precisely, for an
arithmetic function f(x) there exists an asymptotic if
one has an asymptotic indentity
j(x) = <j>(x)+R(x),
where cp(x) is the approximating function, and R(x) is
the error term, about which one knows in general only
that
lim..BQl = o.
x~oo<j>(x)
A short notation is: f(x)=cp(x)+o(cp(x or f(x)~cp(x)
(cf. Asymptotic formula).
Finding asymptotics of arithmetic functions is one of
the most important problems in analytic number
theory. This is explained by the fact that almost-all
arithmetic functions with interesting arithmetical pro-
perties are characterized by extreme irregularity in their
changes as the argument increases. If one considers
instead of an arithmetic function f(x) its average value
(~n .;;)"(n / x (n a natural number), then the 'irregu-
larity' of f(x) is smoothed out. Hence, a typical prob-
lem for an arithmetic function is that of obtaining an
asymptotic for its average value function. For example,
the average value of the function T(n), whose value at n
is the number of divisors of n, is equal to
I
- ~ T(n) ~ Inn.
n m~n
The problem that arises here, of the best possible
bound for the error term in the asymptotic identity, is
still unsolved (1984) for many functions, in particular
for the function T(X) (cf. Analytic number theory).
Asymptotics of arithmetic functions play an impor-
tant a role in additive problems (cf. Additive number
theory). For many of them there is no known direct
proof of the existence of decompositions of a number
into terms of a given form. However, as soon as one
has an asymptotic for the number of decompositions of
the type one is looking for, one can deduce that the
required decomposition exists for all sufficiently large
numbers n.
B.M. Bredikhin
AMS 1980 Subject Classification: 10H25
ATOM - A minimal non-zero element of a partially
ordered set with a zero 0, i.e. an element p such that
275
ATOM
O<x~ implies x=p.
L.A. Skornyakov
Editorial comments. By a natural extension of meaning,
the term atom is also used for an object of a category hav-
ing no subobjects other than itself and the null subobject
(cf. Null object of a category).
AMS 1980 Subject Classification: 06A 10
ATOMIC DISTRIBUTION - A probability distribu-
tion on a a-algebra s1 of subsets of a space Q concen-
trated on the set of its atoms (cf. Atom). A set A E s1
in a probability space (Q, s1, P) is called an atom of the
distribution if peA 0 and if it follows from A ~ BE s1
that P(B)=O or that peA \B)=O. In the particular
case in which the atoms are discrete points, an atomic
distribution coincides with a discrete distribution.
A. V. Prokhorov
AMS 1980 Subject Classification: 60EXX
ATOMIC LATTICE - A lattice with a zero in which,
for each non-zero element a, there exists an atom p ";;a.
o.A. Ivanova
AMS 1980 Subject Classification: 06BXX
ATOMIC RING - An integral domain with a unit ele-
ment, satisfying the ascending chain condition (max-
imum condition) for principal ideals. In other words,
any family of principal ideals of an atomic ring has a
maximal element. An element of the ring is called an
atom, or an extremal element, or an indecomposable ele-
ment, if it cannot be decomposed into a product of
non-invertible elements. An integral domain is an
atomic ring if and only if each non-zero non-invertible
element is a product of atoms. All Noetherian rings are
atomic. Factorial rings are atomic rings in which each
atom is a prime element (i.e. an element that generates a
prime ideal). An atomic Bezout ring is a principal
ideal ring. v.I.
Editorial comments. There seems to be no special name
established in Western literature for these rings. Instead one
simply speaks of rings which satisfy the ascending chain
condition for principal ideals or the divisor chain condition.
References
[A1) GILMORE, R.: Multiplicative ideal theory, M. Dekker, 1972.
AMS 1980 Subject Classification: 13EXX,
1 6A02, 16A 14, 16A 1 5, 16A34
ATTAINABLE BOUNDARY ARC of a domain G in
the z-plane - A Jordan arc forming part of the boun-
dary of G and at the same time forming part of the
boundary of some Jordan domain g C G. Each point
on an attainable boundary arc is an attainable (from
the inside of g) boundary point of G (d. Attainable
276
boundary point). A conformal mapping of a simply-
connected domain G onto the unit disc D: I z I < 1 can
be continuously extended to the non-terminal points of
an attainable boundary arc to a homeomorphism of the
open attainable boundary arc onto some open arc of
the circle I z I = 1. E.P. Dolzhenko
Editorial comments. The standard Western terminology
is accessible boundary arc and accessible boundary point,
see e.g. [A 1].
References
[A1) COLLINGWOOD, E.F. and LOHWATER, A.I.: The theory of
cluster sets, Cambridge Univ. Press, 1966.
[A2] MA1ucUSHEVICH, A.I.: Theory of functions of a complex vari-
able, 3, Chelsea, 1977, Chap!. 2 (translated from the Russian).
AMS 1980 Subject Classification: 30040
ATTAINABLE BOUNDARY POINT - A point on the
boundary of a domain together with the class of
equivalent paths leading from the interior of the
domain to that point. Let ~ be a point on the boundary
aG of a domain G in the complex z-plane and let there
exist a path described by the equation z =z(t), where
the function z(t) is defined and continuous on a certain
segment [a, ,8], Z(t)EG if a";;t<,8, z(,8)=~. One then
says that this path leads to the point ~ (from the inside
of G) and defines the attainable boundary point
represented by f Two paths leading to ~ are said to be
equivalent (or, defining the same attainable boundary
point) if there exists a third path which also leads to ~
from the inside of G and which has non-empty intersec-
tions inside G as close to ~ as one pleases with each of
the two paths considered. The totality of a point ~EaG
and the class of equivalent paths leading to ~ from the
interior of G is said to be an attainable boundary point
of the domain G. Not every point ~EaG represents an
attainable boundary point; on the other hand, the same
point ~EaG can represent several, or even an infinite
Danilov set of different, attainable boundary points.
An attainable boundary point is the unique point of
a prime end (cf. Limit element) of the first kind; a
(multi-point) prime end of the second kind contains
exactly one attainable boundary point, while prime
ends of the third and fourth kinds do not contain
attainable boundary points. Each point of the boun-
dary of a Jordan domain is attainable.
13FXX,
References
[I) MARKUSHEVICH, A.I.: Theory ojjunclIons of a complex variable.
3, Chelsea, 1977 (translated from the Russian).
(2) COLLINGWOOD, E.F. and LOHWATER, A.J.: The theor}' ofclus.
ter sets, Cambridge Univ. Press, 1966.
E.P. Dolzhenko
Editorial comments. The standard Western terminology
is accessible boundary point.
AMS 1980 Subject Classification: 30040
 ATTRACTION. PARTIAL DOMAIN 01

ATIAINABLE SUBGROUP - A subgroup H that can F(-x) c,

- ~ - - as x~oo,
be included in a finite normal series of a group G, i.e. [l-F(x)+F(-x)] C,+C2
in a series for some Cj, C2~0, Cl +C2>0, determined by Vex),
{I} c H = Hoe H, c ... C H. =G and [1- F(x)+ F( - x)] a
[1-F(tx)+F(-tx)] ~t asx~oo,
in which each subgroup Hi is a normal subgroup in
Hi + 1. The property of a subgroup to be attainable is for each constant t >0. Restriction on the behaviour of
transitive. An intersection of attainable subgroups is an the normalizing coefficients Bn, n = 1, 2, ... , leads to
attainable subgroup. The subgroup generated by two narrower classes of distribution functions for which the
attainable subgroups need not be an attainable sub- convergence in distribution (*) holds. The set of distri-
group. A group G all subgroups of which are attainable bution functions F(x) for which (*) converges weakly,
satisfies the normalizer condition, i.e. all subgroups for a suitable choice of An, C >0 and Bn =cn -1 / 2 ,
differ from their normalizers (cf. Nonnalizer of a sub- n = 1,2, ... , to a stable distribution function Vex)
set). Such a group is therefore locally nilpotent. with exponent a, is called the normal domain of attrac-
tion for V(x). The normal domain of attraction of a
References
[I] KUROSH, A.G.: The theory of groups, 1-2, Chelsea, reprint, normal distribution coincides with the set of non-
1955-1956 (translated from the Russian). VM v t degenerate distributions with a finite variance.
. . I\.0PY ov
The normal domain of attraction of a non-degenerate
Editorial comments. Instead of attainable subgroup, the stable distribution function V(x) with exponent 0:
term accessible subgroup is used in [1]. In the Western (0<0:<2) is formed by the functions F(x) for which
literature the term subnormal subgroup is standard for this
kind of subgroup. lim
x~-oo
IF(XI~
x
= c, ;;;.. 0,
References
lim 1 - F(x) -- ::;;, 0
C2 ".-
[A1] SUZUKI, M.: Group theory, 2, Springer, 1986. x--+oo xa
AMS 1980 Subject Classification: 20F26 exist and are finite, where Cj, C2 are determined by
Vex).
ATfRACflON DOMAIN OF A STABLE DISTRIBU-
References
TION, domain of attraction of a stable distribution - The [I] GNEDENKO, B.Y. and KOLMOGOROV, A.N.: Limit distributions
totality of all distribution functions F(x) such that for for sums of independent random variables, Addison-Wesley,
a sequence of independent identically-distributed ran- 1954 (translated from the Russian).
[2] IBRAGIMov, LA. and LINNIK, Yu.V.: Independent and station-
dom variables Xl, X 2 , . , with distribution function ary sequences of random variables, Wolters-Noordhoff, 1971
F(x) and for a suitable choice of constants An and (translated from the Russian).
En >0, n = 1, 2, ... , the distribution of the random [3] PETRov, V.V.: Sums of independent random variables, Springer,
variable 1975 (translated from the Russian).
B.A. Rogozin

~Xk-A. AMS 1980 Subject Classification: 60E07, 60F05
-,,-k~---,''--_ _ , n =1,2, ... , (*)
B.
ATfRACflON, PARTIAL DOMAIN OF, of an
converges weakly, as n-Hx:;, to a non-degenerate distri- infinitely-divisible distribution - The set of all distribu-
bution function V(x), which is necessarily stable. tion functions F(x) such that for a sequence of
One of the fundamental problems in the theory of independent identically-distributed random variables
stable laws is the description of domains of attraction XI, X 2, . . . , with distribution function F, for an
of stable laws. Thus, for the normal distribution, A.Ya. appropriate choice of constants An and En >0,
Khinchin, W. Feller and P. Levy established in 1935 n = 1,2, ... , and a subsequence of integers
that F(x) belongs to the domain of attraction of a nor- n 1 <n2 < .. " the distribution functions of the ran-
mal law if and only if, as X-HlO, dom variables

x2 IY~XdF(y)
!
~o.
y2 dF(y)
Iy <x
converge weakly, as k-HIJ, to a (given) non-degenerate
Later B.V. Gnedenko and W. Doeblin (1940) gave a distribution function V(x) that is infinitely divisible;
description of the domain of attraction of a stable law every infinitely-divisible distribution has a non-empty
with exponent a, 0<a<2: F(x) belongs to the domain partial domain of attraction. There exist distribution
of attraction of a non-degenerate stable law V(x) with functions that do not belong to any partial domain of
exponent 0: if and only if: attraction and there also exist distribution functions

277
ATTRACTION. PARTIAL DOMAIN OF
that belong to the partial domain of attraction of any
infinitely-divisible distribution.
References
[1) GNEDENKO, B.V. and KOLMOGOROV, A.N.: Limit distributions
for sums of independent random variables, Addison-Wesley,
1954 (translated from the Russian).
B.A. Rogozin
Editorial comments. The notion defined in this article
is also commonly called the domain of partial attraction.
AMS 1980 Subject Classification: 60E07, 60F05
AUTO-CORRELATION of a stochastic process XI -
Correlation of the values of XI and XI +h. The term
'auto-correlation', along with the term 'correlation func-
tion', is mostly employed in studies of stationary sto-
chastic processes, in which the auto-correlation depends
only on h and not on t (d. Stationary stochastic pro-
cess).
A. V Prokhorov
Editorial comments. I.e. the auto-correlation of the pro-
cess Xt is the correlation coefficient of Xt and Xt + h.
AMS 1980 Subject Classification: 62M10
AUTO-CORRELOGRAM of a stochastic process - The
same as the correlogram.
AMS 1980 Subject Classification: 62M10
AUTO-OSCILLATION - A non-damped oscillation
in a non-linear dynamical system, whose amplitude and
frequency can remain constant during a long period of
time, are largely independent of the initial conditions,
and are determined by the properties of the system
itself (see also Oscillations, theory of). The term 'auto-
oscillation' was introduced by A.A. Andronov (see [1 D.
Dynamical systems capable of performing oscilla-
tions are known as auto-oscillating systems. They
include clocks, generators of electric vibrations, electric
bells, blow and bow musical instruments, etc. Under
certain conditions auto-oscillations can also be pro-
duced in dynamical systems which normally function
without auto-oscillations (e.g. auto-oscillations in the
front suspension of a car-'noise', the flutter of an air-
craft wing, or auto-oscillations in automatic control and
monitoring systems). The simplest auto-oscillating sys-
tem can be represented as consisting of a constant
source of energy, a mechanism regulating the supply of
energy to the oscillating system, and the oscillating sys-
tem itself. An essential feature of such a system is its
feedback nature: on the one hand, the regulating
mechanism controls the motion of the oscillating sys-
tem but, on the other, it is the motion of the oscillating
system that influences the operation of the regulating
mechanism [6]. From the mathematical point of view,
autonomous auto-oscillating systems with one degree of
27R
freedom may be defined as systems whose equations of
motion have one or more limit cycles in the phase
plane (cf. Limit cycle).
An important typical property of auto-oscillations is
that their amplitude is largely independent of the initial
conditions, i.e. the existence of one or more ranges of
initial conditions such that to any initial set of condi-
tions within anyone of these ranges there corresponds
the same auto-oscillation amplitude. This means that
oscillating systems may contain several stationary
processes with different amplitudes; each one of these
becomes realized in the system, depending on a choice
of a region of initial conditions. This property consti-
tutes the basic difference between periodic motion in
an auto-oscillating system and periodic motion in a
conservative system. Another typical property of auto-
oscillating systems is the fact that the periods of the
auto-oscillations are determined by the properties of
the system itself and are not imposed from outside.
This is the basic difference between auto-oscillations
and forced oscillations.
An essential feature of auto-oscillations consists in
the fact that the loss of energy must be compensated by
a constant source of energy. Since in an autonomous
system the forces do not explicitly depend on time, this
source of energy must produce a force which is not a
given function of time and which is determined by the
system itself. Constant sources of energy include the
winding mechanism in a clock, a shaft rotating at a
constant angular velocity, an endless belt moving at a
constant rate, or a jet of liquid or gas having a constant
velocity; in systems in which the motive force is an
electric flow, an electric battery may serve as such a
source.
An example of a very simple auto-oscillating system
is a pendulum in a viscous frictional medium, acted
upon by a force of constant magnitude which always
acts in the direction of the motion. The differential
equation of this dynamical system is
x+2hx+k 2 x = F sign x,
where h >0, F>O, k 2 are constant coefficients. Under
any initial conditions a periodic oscillation becomes
established in the system, the maximum deflection of
the pendulum from its position of equilibrium (ampli-
tude) being
] +e- hT F
XII = ]-e hT k2 '
where T=1T / Yk 2 -h 2 is the half-period of the
periodic motion. A stable limit cycle in the phase plane
xx corresponds to this motion [3].
A typical property of all auto-oscillating systems is
the connection between the constant source of energy
and the system. which is such that the energy given by

the source varies periodically, the period of this varia-
tion being determined by the properties of the system.
This may be expressed by saying that an auto-
oscillating system is a system which executes a periodic
process at the expense of a non-periodic energy source.
The shape of auto-oscillations can be close to
sinusoidal, but can also be very different. The latter
type of auto-oscillations is known as relaxation oscilla-
tion. Auto-oscillations with a shape close to sinusoidal
usually occur in systems in which the loss in energy
during one period is small, and consequently, the
amount of the energy fed into the system is also small.
They may also occur when the energy losses in the sys-
tem during one period are large, but if the parameters
are properly chosen, the losses are compensated not
merely for each period, but also for each small fraction
of a period. Such systems include the so-called Re-
generators of sinusoidal oscillations [4]. During relaxa-
tion oscillations energy losses are usually large and are
compensated for after one period by almost the entire
energy of the oscillating system.
The generation of auto-oscillations may be 'soft' or
'hard' [2]. In the former case, variation of some param-
eter of a system in a stable equilibrium results in the
generation of auto-oscillations whose amplitude steadily
increases from zero as the parameter is continuously
varied. If the parameter is now varied in the opposite
direction, the amplitude of the auto-oscillation steadily
decreases to zero, and the system returns to a stable
equilibrium. If the generation of the auto-oscillations is
hard, the system passes from a state of stable equili-
brium to an auto-oscillation with finite amplitude as
one of its parameters is slowly and continuously varied;
as the parameter is further varied, the rate of increase
of the amplitude becomes constant. If the parameter is
varied in the reverse direction, the amplitude continu-
ously decreases and, after it has attained a certain
value, the system returns to the state of stable equili-
brium. It is important to note in this connection that
the system passes from stable equilibrium to auto-
oscillations and from auto-oscillations back to stable
equilibrium at different values of the parameter.
Auto-oscillating systems display an interesting and
important property - the effect of enforced synchroni-
zation, which is sometimes called 'capture' or 'mode
locking'. This takes place when the difference between
the frequency of the auto-oscillating system and that of
the external force acting on it is sufficiently small. The
steady periodic motion of the system then assumes the
frequency of the external force, i.e. the external force
imparts its own frequency to the auto-oscillating system
[3], [5].
References
[1) ANDRONOV, A.A: Collected work, Moscow, 1956 (in Russian).
AUTO-REGRESSIVE PROCESS
[2) ANDRONOV, A.A., VITI, A.A. and KHAiKIN, S.E.: Theory of
oscillations, Moscow, 1959 (in Russian).
[3) BUfENIN, N.V.: Elements of the theory of non-linear oscillations,
Blaisdell, 1965 (translated from the Russian).
[4) GoRELIK, G.S.: Oscillations and waves, Moscow-Leningrad,
1950 (in Russian).
[5) FEL'DBAUM, A.A: Introduction to the theory of non-linear
chains, Moscow-Leningrad, 1948 (in Russian).
[6) KIiARKEvICH, AA: Auto-oscillation, Moscow, 1953 (in Rus-
sian). N. V Butenin
Editorial comments. The general mathematical definition
of 'autonomous auto-oscillating systems' would be: vector
fields having attractors which are non-constant periodic
solutions. The 'ranges of initial conditions' are the basins of
attraction of these solutions. In conservative systems,
periodic solutions cannot be attractors and, if all solutions in
an open set are periodic, the amplitude and period will usu-
ally depend on the initial conditions. The 'soft auto-
oscillations' can be identified with the Hopf bifurcation in
dynamical systems.
In [A1], especially pages 517 and 504 are useful.
References
[A1] ABRAHAM, R. and MARSDEN, J.E.: Foundations of mechanics,
Benjamin-Cummings, 1978.
AMS 1980 Subject Classification: 58F21, 58FXX,
93CXX
AUTO-REGRESSION - A regressive dependence of
the values of Xn of a given random sequence {Xn:
n =0, + 1, ... } on the preceding values of
Xn - I , . ,Xn - m . A linear auto-regression scheme of
order m is defined by a linear regression equation
between Xn and Xn - k , k=I, ... ,m, i.e.
Xn == !31Xn- 1+ ... +!3m Xn-m+n, (*)
where 131, ... ,13m are constants and the random vari-
ables (n are identically distributed with average zero,
variance r? and are uncorrelated (sometimes they are
assumed to be independent). An auto-regression
scheme is a useful stochastic model for the description
of certain time series (the concept of a linear auto-
regression scheme was first introduced by G. Yule in
1921) in order to analyze time series describing a sys-
tem which is oscillating under the effect of internal
forces and random external shocks. The auto-regression
scheme (*) may be regarded as a stochastic process of a
special type: an auto-regressive process.
A. V Prokhorov
AMS 1980 Subject Classification: 62JXX
AUTO-REGRESSIVE PROCESS - A stochastic pro-
cess {Xt } whose values satisfy an auto-regression equa-
tion with certain constants aI, ... , ap :
X, == a1X'-1 + ... +apX,_p + Y" (*)
where p is some positive number and where the vari-
ables Yt are usually assumed to be uncorrelated and
279
AUTO-REGRESSIVE PROCESS

identically distributed around their average value 0 AUTOMATA, ALGEBRAIC THEORY OF - A branch
with a variance 01. If all the zeros of the function of the theory of automata (cf. Automata, theory of) in
cp(z) =zP - a 1zP -I - . . . - ap of a complex variable z which algebraic tools are employed in the study of
lie inside the unit circle, then equation (*) has the solu- automata. The algebraic theory of automata is based on
tion 00
the fact that automata can be regarded as special alge-
~ = ~b.Y,-" bras or algebraic systems. In addition, events which are
.=0
representable by finite automata form an algebra with
where the b. are connected with the aj by the relation respect to the operations of union, composition and
iteration, generated by a finite set of so-called elemen-
tary events, each one of which consists of a one-letter
For example, let {Yt } be a white noise process with or empty word. The algebraic approach makes it possi-
spectral density 01 /2'IT; in such a case the only kind of ble to utilize algebraic results in the theory of automata
auto-regressive process satisfying equation (*) will be a directly, and in certain cases aids in establishing a con-
process {Xt } with spectral density nection between the theory of automata and other
branches of mathematics. E.g., the theory of automata
f(A) = [~) 1 cp(e iA ) 1- 2 , was used to obtain a proof of the solvability of certain
second-order arithmetical theories, and also to obtain a
which is stationary in the wide sense if </l(eo,) has no new, simpler solution of the Burnside problem.
real zeros. The autocovariances (cf. Autocovariance) from the algebraic point of view, a finite or infinite
rk = EXtXt -k of the process satisfy the recurrence rela- automata 2f=(A, S, B, cp,l[;) is a tribasic algebra, i.e. an
tion algebra with three sets of elements A, S and B and two
rk = al'k-l + ... +ap'k-p, k>O,
operations cp: S XA ~S and 1[;: S XA ~B. On the other
and, in terms of the b., have the form hand, a transition system (A, S, cp), where
00
A ={ai, ... ,an} is the input alphabet, S is the set of
'k = rl- ~b.b.+k' states (cf. Automaton, fmite) can be regarded as an
.=0
The parameters ak of the auto-regression are connected algebra <S, a b . . . ,an> with unary operations,
with the auto-correlation coefficients Pk =rk / ro of the denoted by the letters aj of the input alphabet, such
process by the matrix relation that ajsj = CP(Sj , aj). Concepts such as a homomorphism
of automata (cf. Automata, homomorp1sm of), an iso-
A = R-lp*
P P' morphism, a sub-automaton, etc., as applied to auto-
where A =(ab ... ,apl, P;=(PI,'" ,pp ) and Rp is mata, are thus naturally defined. This approach also
the matrix of auto-correlation coefficients (the permits the assignment to the automaton
Yule - Walker equation). 2f=(A, S, E, cp,l[;) of a semi-group of transformations
References of the set S (under the composition operation), gen-
[l] GRENANDER, U. and ROSENBLAIT, M.: Statistical analysis of erated by the operations aj, so that for any arbitrary aj,
stationary time series, Wiley, 1957.
[2] HANNAN, EJ.: Time series analysis, Methuen, London, 1960.
aj from A and S from S one puts
A. V. Prokhorov ajajS = cp(cp(s, aj), aj).
AMS 1980 Subject Classification: 62JXX
This semi-group is known as the semi-group of the auto-
AUTOCOVARIANCE of a stochastic process Xl - The maton 2 and is used as a means of describing certain
covariance of Xl and Xl + h . If EX denotes the properties of automata (classification, decomposability,
mathematical expectation of a random variable X, then isomorphism, etc.) in algebraic language. Also, to any
the autocovariance equals semi-group Q with a unit element there may be
assigned an automaton with a given input alphabet
E(Xt - EXt )(Xt +h -- EXt f h)'
A = { ai, ... ,an} and a set of states Q as follows. To
The term 'autocovariance' is usually applied to station- each letter aj of A is assigned some element qj from Q,
ary (in the wide sense) stochastic processes (d. Station- and the transition function is then defined as
ary stochastic process). For such processes the autoco- </l(q, aj)=qqj. The semi-group of such an automaton is
variance depends only on h and differs from the auto- isomorphic to the sub-semi-group of Q generated by
correlation by the presence of a single factor which is the elements q I, . . . ,qn, which means that if
equal to the variance of XI' The terms 'covariance q 1, . ,qn are generators of Q, then the semi-group of
function' and 'autocovariance function' are used the automaton is isomorphic to the original semi-group
together with the term 'autocovariance'. Q. The semi-group of an automaton is clearly iso-
A. V. Prokhorov morphic to a quotient semi-group of the input semi-
AMS 1980 Subject Classification: 62M10 group A' of all words over the alphabet A with the

280

operation of successive union (concatenation) of words
according to the congruence
R = {(a, fJ): a, fJEA' and for all s ES(I[>(S, a)=l[>(s, fJ))}.
For any state s ES the congruence R is the maximal
subcongruence of the relation
Rs = {(a, fJ): a, fJEA' and I[>(s, a)=l[>(s, fJ)}.
This means, in particular, that an event which can be
represented by the initial acceptor 2(s =(A, S, cp, S') is
a union of certain R-classes. Since the semi-group of an
automaton characterizes it up to an isomorphism, vari-
ous classes of semi-groups have their own classes of
automata. If the semi-group of an automaton is free,
Abelian, cyclic, nilpotent, etc., or is itself a group, the
automaton is known, respectively, as a free, Abelian,
cyclic, nilpotent or a group automaton; these last are also
called permutation automata. Another approach, con-
nected with the algebraic classification of transition and
output functions, leads to the classes of linear, subsitu-
tion and other automata (cf. Automaton). Substitution
automata realize one-to-one functions and are used in
coding theory. Linear automata are of interest because
of their simple structure.
An automaton (A, S, B, CP,1/I) is said to be a linear
automaton if A, S and B are linear spaces over some
field P,
I[>(s, a) = (l>I(s)+~(a), #.s, a) = lh(s)+1jJ2(a),
where CPI, <f>2, 1/11, 1/12 are, respectively, linear mappings of
S into S, A into S, S into B, and A into B. It is usually
assumed that the field P is finite, and that the spaces
A, S and B are finite-dimensional; in such a case a
linear automaton is a finite automaton. If multiple
operations are permitted in the representation of a fin-
ite acceptor as an algebra whose operations are letters
of the input alphabet, the resulting generalization is
known as an automaton over terms (automaton over
trees; generalized automaton). Such automata are used
to prove the solvability of certain second-order
mathematical theories.
References
[1) Algebraic theory of automata, languages and semi-groups, Mos-
cow, 1975 (in Russian).
[2] GLUSHKOV, Y.M.: 'The abstract theory of automata', Russian
Math. Surveys 16, no. 5 (1961), 1-53. (Uspekhi Mat. Nauk 16,
no. 5 (1961), 3-62)
(3) THATCHER, J.W. and WRIGHT, J.B.: Kibernet. Sb. 6 (1969),
112-114.
See also the references to Automaton.
V.B. Kudryavtsev
Yu.I. Yanov
AMS 1980 Subject Classification: 18820, 68D05
AUTOMATA, COMPLETE SYSTEMS OF - Special
subsets of a given class M of automata on which a cer-
AUTOMATA. COMPLETE SYSTEMS OF
tain set 0 of operations with values in M has been
defined. These subsets have the following fundamental
completeness property: The set of all automata obtained
by a finite number of operations from 0 to automata
from a given subset A c;;,M, is identical with M. The
problem of whether a set is complete or not is known
as the completeness problem for automata. It has been
studied for various models of automata and operations.
The range of objects of such studies includes, in order
of increasing complexity, true automata, automata real-
izing functions with delays (i.e. functions of k-valued
logic with a certain shift in time), finite automata, etc.
The completeness problem of true automata with super-
position operations (which are usually considered) is
essentially the same as the completeness problem for
functions of k-valued logic (see Many-valued logic) and
has been fairly thoroughly studied. Synchronous super-
position is defined as superposition of automata for
which the automata realizing functions with the same
delay are combined at all inputs. The criteria of com-
pleteness, found under these conditions, imply, in par-
ticular, the existence of an algorithm by means of
which the completeness or incompleteness of any finite
system of automata can be established. Fundamental
criteria of completeness are given in terms of the so-
called pre-complete classes (i.e. subsets of the class M
that are closed with respect to the operations from 0,
each one of which forms together with any automaton
not belonging to it a complete system). It has been
shown that a set A is complete if and only if it is not a
subset of any pre-complete class. These classes, in tum,
have been fully described. This approach was success-
fully employed in a large number of other cases. It is
also applicable, in principle, to finite automata if the
operation of superposition of automata and the feed-
back operation are chosen as 0 (cf. Automata, composi-
tion of). The criterion given above is applicable also in
this case; it has been established, however, that the
family of pre-complete classes is a continual class,
which excludes the existence of any effective criteria in
these terms. One may note that in all the cases men-
tioned here there exist finite complete systems, and for
this reason the completeness problem for arbitrary sys-
tems of automata is reduced to the completeness prob-
lem for finite systems of automata. One must consider
in this context, as above, the problem of algorithmic
solvability of the completeness problem for finite sys-
tems of finite automata. The problem may be general-
ized as follows: To determine, for a given automaton a
and set B, whether or not a may be obtained from
automata in B by effecting a given set of operations.
One must thus study the predicate P (x, y) - 'the auto-
maton x is realized by means of a set y'. It was found
that the problem of recognition of realizability is algo-
281
AUTOMATA, COMPLETE SYSTEMS OF

rithmically unsolvable for any given a, i.e. that the ;=1, ... ,n, is the automaton 2l=(A, S, B, Cp,o/) for
one-place predicate P (a, y) is non-recursive. On the which
other hand, for various values B of the parameter y, the S = SIX ... XSn , A =AIX ... XA n ,
predicate P(x, B) may be either recursive or non-
recursive. A problem connected with the algorithmic B = BI X .. , XBn
unsolvability of the completeness problem for automata while the functions <p(s, a) and o/{s, a) are defined by
is the problem of finding classes of sets for which it has the relations
an effective solution. In particular, there exists an algo-
rithm for the recognition of the completeness of sys-
tems consisting only of Moore automata and all true = (4)I(SIo al), ... ,4>n(sn, an,
automata. A related problem is that of finding concrete 1/{(SIo ... ,sn)' (a 10 ... ,an =
complete sets of automata with given properties. It was = (lh (s 10 al), ... ,I/In(sn, an.
found that, for any positive integer n, there exists a
The superposition of automata
complete system of automata no proper subset of which
211 =(A h SI, B I , </>1, 0/1) and 212 =(Bh S2, B 2 , </>2, 0/2) is
is complete, and that, for a given n, the number of such
the automaton ~ = (A I, S, B 2, cp, 0/) for which
systems is infinite. There also exists a most simple, in a
S=S2 XS I and
certain sense, automaton with two states, two input
channels and one output channel that constitutes a 4>(S2,SI),a) = (f/>z(S2, I/I1(SIo a, 4>1 (s 10 a,
complete system. The completeness problem has also 1/{(S2,sl),a) = 1/12(S2,lh(sloa,
been studied for different generalizations of finite auto-
where (S2, SI)ES and a EA I.
mata and operations performed on them; another kind In problems of completeness and synthesis of auto-
of generalization concerns the introduction of various mata the feedback operation is very important. This
equivalence relations into the class of automata under operation is applicable to automata with n inputs and
consideration. m outputs 2l=(A, S, B, </>,1/;) where A =A I X ... XA n ,
References B=BI X ... XB m ,
[I] YABLONSKII, S.V.: 'Functional constructions in k-placed logic',
Trudy Mat. Inst. Steklov. SI (1958), 5-142 (in Russian). 1/{s, (a 10 ... ,an =
[2] KUDRYAVTSEV, V.B.: 'The completeness theorem for one class
of feedback automata', Problemy Kibernet. 8 (1958),91-115 (in
= (I/II(S, (a 10 ,an, ... ,I/Im(s, (a 10 . ,an),
Russian). such that for certain i, j the relation B; =Aj is true and
[3] KUDRYAVTSEV, V.B.: 'On the cardinality of sets of pre-
complete sets of certain function systems, related to automata', the function 0/; is independent of aj, i.e.
Problemy Kibernet. (1965),45-74 (in Russian).
[4] KRATKO, M.l.: 'On the existence of non-recursive bases of fin-
I/Ii = I/Ii(S, (a 10 ... ,aj-I, aj+lo ... ,an,
ite automata', Algebra i Logika 3, no. 2 (1964), 33-44 (in Rus- Under these conditions, the feedback operation on
sian).
[5] LETICHEVSKll, A.A.: Conditional completeness in the class of the ;-th output and j-th input of automaton 21 yields
Moore automata, Kiev, 1963 (in Russian). the automaton 2f = (A', S, B, </>',1/;') such that
[6] BUEVICH, V.A.: 'The construction of universal finitely- A'=AIX'" XA} - IXA+IX'" XA n,
}
determined functions in two input variables, having two inter-
nal states', Problemy Kibernet., no. 22 (1970), 75-83 (in Rus- <I>'(s, (aI, ... ,aj-lo aj+l, ... ,an =
sian).
V.B. Kydryavtsev = 4>(s, (aI, ... ,aj-lo
Editorial commp.nts. The topic of this article is rather
esoteric. The problem of completeness (in this sense) has
1/1' (s, (a 10 , aj - 10 aj + I, ... , an =
not really been studied outside the Soviet Union and is not
considered very significant in the Western literature. = 1/{s, (aI, ... ,aj_1>
AMS 1980 Subject Classification: 68D05, 03D05 "'i(S, (al,. ,aj_l, aj+I, .. ,an, aj+I, .. ,an,

There are also other kinds of compositions of auto-

AUTOMATA, COMPOSmON OF - Operations by
which more complex automata are obtained from other
automata by combining the initial automata in accor-
dance with certain rules. Composition of automata
plays an important role in problems of synthesis and
decomposition of automata. The most important and
the most frequently used compositions of automata are
the direct product, superposition and feedback.
The direct product of automata 21; =(A;, S;, B;, cp;, 1/;;),
mata, e.g. the product, direct sum, semi-direct product,
etc.
References
[1) GLUSHKOV, V.M.: 'The abstract theory of automata', Russian
Math. Surveys 16, no. 5 (1961), I-53. (Uspekhi Mat. Nauk 16,
no. 5 (1961), 3-62)
[2) KUDRYAVTSEV, V.B.: 'On the cardinality of sets of pre-
complete sets of certain function systems, related to automata',
Problemy Kibernet. (1965),45-74 (in Russian). VN. Red'ko

282

Editorial comments. For more information see [A1].
References
[A 1] GEISEG, F.: Products of automata, Springer, 1986.
AMS 1980 Subject Classification: 68005, 20M35
AUTOMATA, EQUIVALENCE OF - An equivalence
relation on the set of automata, which arises in the con-
text of studying some individual internal properties of
automata. Such a property is usually the behaviour of
automata (cf. Automaton, behaviour of an), since two
automata are considered equivalent if their behaviour is
identical. In this context the behaviour of an automa-
ton is usually understood to be the totality of functions
realized by it (cf. Automaton, fmite). For finite auto-
mata such an equivalence relation is solvable, and for
this reason there exists a minimization algorithm for
automata, i.e. an algorithm used to construct for a
given automaton an equivalent automaton with the
smallest possible number of states (a minimal automa-
ton).
Equivalence of automata is conveniently defined in
terms of the concept of equivalence of states: Two states
s and s' of the automata ~ and ~' (they may be identi-
cal) are said to be equivalent if the initialized automata
~s and ~~. realize the same function. Such an
equivalence of the automata ~ and ~' will then mean
that it is possible to find, for an arbitrary state of ~, a
state of ~' equivalent to it, and vice versa. An automa-
ton is minimal if and only if no two states of it are
equivalent. The minimal automaton is unambiguously
defined for any automaton, up to an isomorphism. The
algorithm of solution of this equivalence relation for
finite automata is based on Moore's theorem, according
to which states s and s' of ~ are equivalent precisely if
the functions realized by the initialized automata ~s
and ~s are identical on words of length n -1, where n
is the number of states of ~. If the states sand s'
belong, respectively, to two automata ~ and ~', the
estimate is equal to n +n' -1, where n' is the number
of states of 2{'. This theorem also forms the principle of
a well-known minimization algorithm of finite auto-
mata, which involves the construction of the so-called
reduced automaton, the states of which are classes of
equivalent states, while its transition and output func-
tions are naturally induced by the corresponding func-
tions of the original automaton. A reduced automaton
is minimal, since any two of its states are non-
equivalent.
There exist asymptotic estimates of the number
A (m, n, p) of minimal automata with n states, m input
and p output letters; if m+n+p-HX) for m;;;;;'3, the
estimate is A (m, n, p ),....,,(np )mn / n !, while
(n" \2"
A(2, n,p) ~ e-l/(2P')~.
n.
AUTOMATA. EXPERIMENTS WITH
Another problem in the study of equivalence of auto-
mata is the problem of equivalent transfonnations of
automata. This problem was studied for two of the pos-
sible ways of defining finite automata - diagrams and
logical networks. Generally speaking, the problem con-
sists in finding a system of transformation rules which
satisfy certain conditions and which enable one to
transform any given automaton into an automaton
equivalent to it: a so-called complete system of rules. In
both cases the transformation rule is a pair of schemes
(fragments of diagrams or logical networks) which real-
ize identical choices of mappings. The rule is applied
by substituting one fragment for another. For finite
automata a complete system of rules cannot exist, but
they do exist for logical networks with a limited
number of delays.
The fundamental concepts, problems and methods
which arise in the study of the equivalence of finite
automata are usually extended to other types of auto-
mata, with due regard for their particular features. For
references see Automaton.
V.B. Kudryavtsev
Yu.I. Yanov
Editorial comments. The reader is referred to the arti-
cles Formal languages and L-systems. Most of the topics
in this article are nowadays not considered to be in the
mainstream anymore.
AMS 1980 Subject Classification: 68005
AUTOMATA, EXPERIMENTS WITH - A method of
obtaining information about the internal structure of an
automaton from its behaviour as can be deduced from
external experiments (that is, experiments where input
words are fed into an automaton, the corresponding
sequence of output words is examined and conclusions
are drawn on the basis of these observations).
Experiments with automata can be used to seek
approaches to the solution of the following problems:
1) Given that an automaton is in an initial state it is
required to determine this state.
2) The construction of an experiment that transfers
an automaton from any state to some pre-assigned state
(the adjustment problem).
3) The verification of the correct functioning of an
automaton. It is required to find out by means of an
experiment if a given automaton functions correctly.
4) The diagnostic of an automaton. It is required to
find out not only if an automaton is in good working
order, but also what kind of malfunction occurs.
5) The problem of identifying an automaton from a
given class (deciphering). Knowing that the automaton,
viewed as a 'black box', belongs to a given set, it is
required to find out which element the automaton is.
Let ~ be a class of initialized automata with input
283
AUTOMATA, EXPERIMENTS WITH
alphabet X and output alphabet Y. Next, let X be the
collection of all words in X, let aEX and let A em:.
An experiment with the automaton A is a set [a, A] of all
pairs of the form (x,y) such that x Ca and y is a word
into which A produces x.
The classification of experiments. The following types
of experiments are usually considered:
1. A simple unconditional experiment - a set [a, A],
where a consists of a single word. In the process of
experimenting the single element a is fed into the input
of the automaton aLd the resulting output is some
word from Y.
2. A simple conditional experiment - a set [a, A],
where a consists of a single word X=(X(I), ... ,X(T
and where for any j E {1, ... ,T} the values X(j)
depend on [a',A], a'=(X(I), ... ,X(j-I. Thus, in
the process of experimenting with automata every sub-
sequent letter of the input word depends on the output
word obtained earlier.
3. A multiple unconditional experiment - a generali-
zation of a simple unconditional experiment to the case
when the set [a, A] is such that a consists of more than
one word. It is usually assumed that in the process of
execution of the given experiment there is more than
one specimen of the automaton A, or that the automa-
ton can be reset in its initial state.
4. A multiple conditional experiment - a generaliza-
tion of a simple conditional experiment to the case
when the set [a, A] is such that a consists of more than
one word:
a = {(XI(I), ... ,X1(TI,'" , (Xs(l), ... ,X,(TJ)}.
Moreover, two possibilities can be distinguished here:
a) for any i E {l, ... ,s} and j E {l, ... ,T;} the value
of Xi(j) depends only on [a', A], where a' consists of
the single word (Xi(l), ... ,Xi(j -1; or b) for any
iE{1, ... ,s} and jE{l, ... ,T;} the value of x,(j)
depends on [a", A], where a" is the set consisting of the
. \
coII ectlOns {(X\(l), ... ,X\(tj, ... ,
(Xs(l), ... ,X,(tj} and the
index t;
for i E { 1, ... ,s} coincides with Ti if j - I >Ti
or withj-l ifj-l:OS:::Ti'
Measures of the complexity of experiments. These are
some numerical characteristics that estimate the degree
of complexity of an experiment. For a multiple experi-
ment the usual measures of complexity are the heighl of
the experiment - the number of symbols of the largest
input word, the multiplicity - the number of input
words, and the length - the total number of symbols
in all the input words. It is considered that the measure
of complexity of a simple experiment is its length (the
number of symbols of the input word used).
Some results. A Moore automaton A with r states, m
input and p output symbols is called an (r, m, p)-
284
automaton. The following theorem on the distinguisha-
bility of states by a simple experiment holds: If A is an
(r, m, p )-automaton with all states distinct, then the dis-
tinguishability of any two states can be established by
means of a simple experiment of length r - 1, and the
bound r -1 cannot be lowered (see [4]).
The problem of determining the initial state of an
automaton can be solved by means of a simple uncon-
ditional experiment of length A, where
3,(I+'d/ 6 ,,;;; A< 3,(1+<2)/6 (j,(2- 0 as r_oo).
Adjustment for an automaton with r states of which
k are admissible can always be solved by means of a
simple unconditional experiment of length A, where
A:OS:::(2r-k)(k-l)/2, and this bound cannot be
improved.
As a rule, the above estimates are exact only for a
small number of automata. It has been established (see
[4]) that the problem of determining the initial state for
aUnost-all (r, m, p )-automata with two admissible states
can be solved by means of a simple unconditional
experiment of length A, where A:os:::log", lo&> r +4, and
the adjustment problem for almost-all (r, m, p)-
automata by means of a simple unconditional experi-
ment of length A, where A< 510&> r.
There is no algorithm that would permit the decy-
phering of initialized Mealy automata with unknown
number of states. However, it turns out that a great
part of such automata can nevertheless be decyphered.
Experiments with automata are also used as tools for
the control of the functioning of automata. A unified
approach has been developed to problems of control of
automata, which makes it possible to indicate funda-
mental methods for their solution. An effective algo-
rithm has been constructed for finding all dead-end
multiple experiments for various automata (see [2]).
References
[1] MOORE, E.F.: 'Gedanken experiments on sequential machines',
in C. Shannon and 1. MacCarthy (eds.): Automata studies,
Princeton Univ. Press, 1956, pp. 179-210.
[2] YABLONSKli, S.V.: 'Mathematical logic, theory of algorithms
and theory of sets', Trudy Mat. [nst. Steklov. 133, 263-272 (in
Russian).
[3] TRAKHTENBROT, B.A. and BARZDIN', YA.M.: Finite automata.
Behaviour and synthesiS, North-Holland, 1973 (translated from
the Russian).
[4] HIBBARD, T.N.: 'Lower upper bounds on minimal terminal
state experiments for two classes of sequential machines', J.
Assoc. Comput. Mach. IS, no. 4 (961),601-612.
Kh.A. Madatyan
AMS 1980 Subject Classification: 68050
AUTOMATA, HOMOMORPlDSM OF - A mapping
of the input and output alphabets and of the set of
states of an automaton into analogous sets of a second
automaton that preserves the transition and output
functions. More precisely, a homomorphism of an auto-
 AUTOMATA, METHODS OF SPECIFICATION OF

maton = (A I, S I, B I, <1>1,1/;1) into an automaton

~I automaton (cf. Automaton, behaviour of an). The func-
~2=(A2,S2,B2,<i>2,1/;2) (cf. Automaton, fmite) is a tions <I> and I/; are often defined by means of a table, a
mapping h=(h l ,h 2 ,h 3 ) of the set AIXSIXB I into diagram or a matrix of transitions. The table of transi-
the set A 2 X S 2 X B 2 tions T of the automaton ~ consists of two subtables
T.p=llttll and T",=llttll, i=l, ... ,n,
hl:A I - A 2, h2: SI - S2, h3: BI - B 2,
j = 1, ... ,m, tt ES, tt EB. The functions <I> and I/; are
such that the following equalities are valid for any s defined by </>(Sj, aj)=tt, I/;(Sj, aj)=tt. If all the
from S I and a from A I : columns of T", coincide, the table T defines a Moore
h 2<PI(S, a) = </l2(h2(S), hl(a)), automaton. For instance, let A 0 = {a I, a2},
h3!/JI(S, a) = !/J2(h2(S), hl(a)). SO={SI,S2}, Bo={bJ,b 2 }; the table T (Fig. 1) will
then specify the functions <Po and 1/;0 of some automa-
Initialized automata are subject to the additional
ton ~ (Fig. 2 shows the subtables T.p and T", of 1).
requirement that h maps the initial state to the initial
state. The automata ~I and ~2 are said to be I
0, I 02
homomorphic if there exists a homomorphism h of auto- I
I
mata mapping A I XS I XB I into A2 XS 2 XB 2 If, in s, s,.b 2 : S2.b2 s, s, I S2 S, b2 : b2
- ----1--- ---t--- ---r--
addition, h is one-to-one, h is called an isomorphism, S2 S2.b, I s,.b 2
I
S2 I
I
S, S2 b, I b2
and the automata ~I and ~2 are said to be isomorphic
automata. If the alphabets A I and A 2, and also the '---~v,..-----' ~--------~yr--------~
Fig. 1. Fig. 2.
alphabets B I and B 2, are identical and the mappings
h I and h 3 are the identity mappings, the homomor-
The diagram of (transitions oj) an automaton is an
phism (isomorphism) h is known as a state homomor-
oriented graph G with a one-to-one correspondence
phism (state isomorphism). Input and output homomor-
between the vertices of the graph and the elements of
phisms (isomorphisms) are defined in a similar manner.
S, and between the edges of the graph and certain sets
State-isomorphic automata and state-homomorphic ini-
of pairs of the form (a, b), a EA, b EB. Each vertex of
tialized automata are equivalent (cf. Automata,
G is the starting point of at least one edge, and the set
equivalence of).
~ of all pairs assigned to the edges with the same ver-
The concept of a homomorphism of automata is used
tex as their ongm has the fonn
in the context of problems concerning minimization,
~={(al,bjJ, ... ,(am,b;J}. The functions <I> and I/;
decomposition and completeness of automata, among
are then defined as follows: </>(Sj, aj)=s" I/;(Sj, aj)=bp if
others.
the pair (aj, bp ) is assigned to the edge issuing from
References from Sj and if this edge ends in the vertex s,. It is con-
[II GLUSHKOV, Y.M.: 'The abstract theory of automata', Russian
Math. Surveys 16, no. 5 (1961), 1-53. (Uspekhi Mat. Nauk 16, venient to fonnulate certain properties of automata in
no. 5 (1961), 3-62) the language of diagrams (connectivity of an automa-
A.A. Letichevskil
ton, attainability of states, etc.). Figure 3 shows the
AMS 1980 Subject Classification: 68005, 20M35 transition diagram of the automaton ~.

AUTOMATA, METHODS OF SPECIFICATION OF -

The various ways of describing automata and their
functioning or behaviour. The method of specification
of an automaton will depend on the approach to the
concept of an automaton. If the macro-approach is
adopted (cf. Automaton, fmite), it is the external
behaviour of an automaton which is described. If the
micro-approach is adopted, the definition should
include a description of the elements from which the
automaton is constructed, and the schemes of their
composition. Methods of specifying finite automata are
described below.
Macro-approach. To specify a finite automaton
~= <A, S, B, <1>,1/;> for given alphabets
A = {a I, ... ,am}, S = {s I, ... , Sn } and
B = {b I, ... , bd is tantamount to describing the func-
tions <I> and I/; or to describing the behaviour of this
Fig. 3.
The transition matrix is used to describe the function-
ing of the transition system <A, S, <1 (cf. Automaton,
finite). It is represented by the (n Xn)-matrix
P = II Pij II whose elements are subsets of the alphabet
A (including the empty subset) such that a EPij if and
only if </>(Sj, a)=sj, so that for any i= 1, ... ,n the rela-
tions Pjp n
P j, = 0 if P =l=r and U n = IPij = A are true.
In order to extend the function <l>i to the set S XA
(where A is the set of all words over the alphabet A
including the empty word), one considers the sequence
of powers of the matrix P. Multiplication of P by itself
is carried out by the usual algorithm, using the opera-
tions of concatenation and union of word sets instead of

285
AUTOMATA. METHODS OF SPECIFICATION OF
multiplication and addition. If a EA is a word of
length d and aEP~f> where P~f> is an element of the
matrix pd, one has q,(Si' a)=Sj. Thus, the transition
matrix P of the transition system <Ao, So, 1/10> and
the matrix p2 have, respectively, the forms:
/I:: :~ II,
P =
{alaI> a2a2} {ala2, a 2a d "
p2 = PXP = {a2al> ala2} {a2 a2, alai} .
"
A number of minimization (reduction) algorithms and
algorithms of automaton synthesis are associated with
these methods of specifying automata.
In order to specify the behaviour of an initialized
(not necessarily finite) automaton ~s, (transformer),
one must define the functionj(a)=t/{sJ, a) which maps
A into B* (or A 00 into Boo, where A 00, BOO are the
sets of all the superwords over the alphabets A and B,
respectively).
Fig. 4.
This function may be defined by means of an informa-
tion tree. From each vertex of the information tree there
issue m edges, which are in one-to-one correspondence
with the respective letters of A. To each vertex is
assigned a state of the automaton ~s" while to each
edge is assigned a letter of B in the following way. The
state S I is assigned to the root. If the state Si has been
assigned to some vertex, the letter b = t/{Si' a) is
assigned to the edge corresponding to the letter a from
A, while to the vertex in which this edge terminates the
state Sj=q,(Si' a) is assigned. To each word
a = al ... ar EA * corresponds a unique sequence
P = PI ... Pr of edges of this tree such that PI issues
from the root and Pi + I issues from the vertex of
termination of Pi' The word f3 = f31 ... f3n where f3i is
the letter from B assigned ~o the edge Pj, i = 1, ... ,r, is
identical with the value t/{s I, a). If the function j is
realized by a finite automaton, the corresponding infor-
mation tree may be effectively determined by means of
a finite subtree. Figure 4 shows Lhe subtree of the ini-
tialized automaton ~~, (the left-hand edges, issuing
from the vertices, correspond to the symbol a I, while
the right-hand edges correspond to the symbol a2)'
The behaviour of a finite automaton (acceptor) in
terms of a representable event (super-event) may be
effected with the aid of a regular expression (cf. Regu-
lar event). Such events may also be specified as sets of
words generated (derived) in some formal system (a
286
semi-Thue grammar, etc.). In this case the semi-Thue
system is given as the quadruple T= <A, S,~, w>
where A, S are finite alphabets, A n S = 0, ~ is an
axiom scheme of the form soX and w is the set of
schemes of derivation rules of the form saX~s' X, r~r,
where s, s', rES, aEA and X is a variable which
assumes values from A *. If, moreover, saX~s' X and
saX~S" X belong to w, one has s' = s". The word a EA
is deducible in the system T if there exists a sequence
of words soa=wJ, ... ,Wn such that Wi+1 is obtained
from Wi, i = 1, ... ,n -1, by the application of a cer-
tain rule from w, Wn ES and w does not contain the rule
wn~wn' A grammar generating a regular event has a
similar form. It is given by the quadruple
f=(A, S, SJ, w) where SI from S is an axiom, w is the
set of rules of the form s;~asj or s;~a. The word
a=al ... an is deducible in f if w includes the rules
SI ~als;" s;, ~a2s;" ... ,Sin ~an' Algorithms which
yield the transition matrix of an automaton from for-
mal schemes of this type are known. Thus, an event
represented in the acceptor <Ao, So, 1/10, SI > by the
state S2 may be defined, for example, as the set of
words that are deducible in the semi-Thue system
.r = <Ao,SO,~={SIX},
"'={Sla2X~S2X; slaIX~sIX;
s2aIX~s2X; s2a2X~sIX; sl~sd>.
There exist several other ways of specifying an auto-
maton. For instance, a (necessarily finite) transition
system <A, S, <P> may be specified as an algebra
< S, P> where P = {<Pa: a EA} is the set of unary
operations on S such that <Pa(s)=q,(a, s). Thus, the
transition system <Ao, So, 1/10> may be regarded as
the algebra <So, {<Pa"<Pa,}>, where <Pa,(SI)=SI,
<Pa,(S2)=S2, <Pa,(SI)=S2,_ <Pf,(S2)=SI' ?ne may also
consider the algebra <S, P> where S is the set of
words of the form sa, S ES, aEA and 4>~ {~a: a EA}
~s the set of unary oper~ti~ns on S such that
<Pa(sa)=saa. The algebra <S, P> is defined by the
system of S-generators and the set of defining relations
w={sja;=s;a;: i=1,2, ... }. Such an algebra specifies
a (usually partial) automaton <A, S, <P> such that if
s;aj =s;a; is a relation from w, the relation
q,(Sj, a;)=q,(s;,a;) is true. For example, the transition
system <110, So, *0> may be specified by a system of
So-generators and the set of defining relations
w= {SI =Slal, SI =S2a2, S2 =Sla2, S2 =S2al}. It IS
then assumed that q,(SI' A)=SI, q,(S2' A)=S2.
The behaviour of an automaton may be described in
the language of the logic of one-place predicates. A
choice of the class of formulas defining a finite automa-
ton may be performed in various ways. The descrip-
tion may be incomplete, because the choice defines a
 AUTOMATA, METHODS OF SPECIFICATION OF

class of automata whose the behaviour is identical to I 2

within the accuracy of this description. For instance, 2 2 3 3
MI
I , M2 = 3 2
,
the 'questionnaire' approach involves specifying classes I
2 2 5 5
of automata with the aid of fragments of information
trees, a partial definition of the functions cp and "', etc. 3
4 4 I 0
The above ways of specifying automata may also be MI
I ,
M2=
I 2 3
applied, with suitable modifications, to the behaviour
2 2 5 5
of certain generalized finite automata (non-
deterministic, infinite, etc.; cf. Automaton). Thus, the In order to define a finite automaton over terms
elements of the table T </> of a finite, non-deterministic <A, S, 0, aa> when the alphabets A ={a I, . . . ,am}
automaton may consist of arbitrary subsets of the set S. and S = {S I, . . . ,sn} are known one must first specify
The behaviour of a finite, non-deterministic acceptor is a mapping 0 of the set A into the finite set of non-
described by a regular expression, as in the case of a negative integers such that there exists at least one ele-
deterministic acceptor. Other generalizations of finite ment a EA for which o(a)=O and, secondly, one must
automata are finite probabilistic automata (cf. Automa- define, for each a; EA, i = I, ... ,m, a o(a;)-ary function
ton, probabilistic), automata over terms, mosaic struc- aa, mapping the set [st., =sx ... XS into S. To
tures, etc. each element a EA for which o(a)=O there corresponds
To specify a probabilistic automaton <A, S, D, J-L> an element aa ES, called an initial state of the automa-
on known alphabets A = {a I, . . . ,am}, ton. For example, if A={al,a2}, S={S],S2},
D={b], ... ,bd, S={SI' ... ,sn}, means to specify, o(ad=O, 0(a2)=2, aal =s], aa,(s], SI)=S2,
for any given i and j (1:OS;;;i:OS;;;n, I:OS;;;j:OS;;;m) an arbitrary aa, (s I, S2) =s], aa, (S2' SI )=S2, aa, (S2' S2)=S2, then the
probability measure J-LiSn bq ) on the set of all pairs functions 0, aa l , aa, define some automaton over terms
(s"b q ), r=I, ... ,n, q=I, ... ,k. To do this, one usu-
<A, S, 0, {aa l , aa,} > with initial state S I.
ally considers a system of square matrices with non-
In order to specify a mosaic structure (an infinite
negative elements
union of transition systems of the form <A, S, cp>,
= II p.r.} II,
where A = S', cf. Automaton), one must define for each
Mq,j

i,r = I, ... ,n; j = I, ... ,m; q = I, ... ,k, integral point of the n-dimensional space an ordered set
such that each matrix M J = "'2:/ = I Mq,J is stochastic. of integral points; this set is called its neighbourhood.
The measure J-L is defined by the relation The input alphabet A of the transition system
J-Lij(sn bq )=J-L7.'/. The probabilistic automaton <A, S, cp> placed at a given point is the Cartesian
<A, S, D, J-L> is considered in conjunction with a cer- product of the set of states of the transition systems
tain initial probability distribution on the set S: placed at the points of its neighbourhood. For exam-
n ple, let S ={S ], S 2 }, and let the ordered set
J.I<J = %(s I), ... , J.I<J(sn, J.I<J(Si);;;' 0, ~J.I<J(Si) = 1. {(a-I, P), (a, P+ I), (0.+ 1, P)} be the neighbourhood
Da,p of an integral point (a, P) of the two-dimensional
i=1

Sometimes probabilistic automata are defined by speci- space. In order to specify a uniform two-dimensional
fying only the matrices MJ or only the matrices mosaic structure, the function cp is defined as follows:
MJ = II p)"q II, where cp(s],sl,sd=sl; in all other cases cp(s',s",Sm)=S2'
Here, the input alphabet A is the Cartesian product
r=1 SXSXS.
i=I, ... ,n; q=I, ... ,k;j=I, ... ,m. Any finite
Micro-approach. A micro-approach definition of a
Markov chain can be regarded as a finite probabilistic
structural automaton describes its constituent elements
automaton in which the matrices MJ, j = 1, ... ,m, are
and the scheme of their combination. This description
identical. The system of matrices shown below defines
may be detailed or more general. It is frequently res-
some probabilistic automaton ~ (n = 2, k = 2, m = 2)
tricted to the so-called canonical scheme of the con-
and the matrices MI, M2, MI, M2 of this automaton:
struction of an automaton. The elements then divide
I 0 I 2 into two groups: functional elements (one-state auto-
4 3 3
MI,I = I I ' M I ,2 = I I '
mata) and memory elements. The canonical scheme
4 4 5 5 (Fig. 5) consists of two functional blocks f and g to
which are added Moore automata
~; = <R;, S;, Q;, cp;, "'; >, i = I, ... ,I, as memory ele-
= o 0
4 2
I
, M 2,2
2 I ' ments. The blocks f and g consist of functional ele-
4 4 5 5 ments. When the structure is defined in this way, these

287
AUTOMATA, METHODS OF SPECIFICATION OF

blocks are not further described, but a specification is where 1 = 1,2, ... , is an integer parameter, So ES,
given (e.g. by means of tables) of the vector functions while the functions j;, i = 1, ... ,I, g;, i = 1, ... ,m, and
realized by these blocks: the variables Xn r = 1, ... ,n, take their values from the
set A. A canonical scheme in which all the memory ele-
ments are identical corresponds to this system:
21'; = <A, S, B, </>,1/;, so>, where A =S =B, i = 1, ... ,I;
</>(a, a')=a', 1/;{a, a')=a. The functioning of the auto-
g = (g1(XI> ... ,Xn' Ul> ... , U/), ... , maton 21'; may essentially be described as follows. At
the moment 1 suppose that the letter a (I) is assigned to
, ... ,g/(XI>'" ,Xn, UI>'" ,U/:
the input of 21';, then the same letter is produced as the
AnXQI X ... XQ/ ~ Hm, output of 21'; at the moment 1 + 1. Let Figure 6
where A n and Bm
are the input and the output alpha- represent an automaton 21' whose canonical equations
bets, respectively, of the canonical scheme. have the form
u(l) = 0,
u(t+l) = x(t)+u(t) (mod 2),
yet) = x(t)Vu(t).

In the general case, a description of structural auto-

mata involves defining a selection of elementary auto-
mata and a certain class of 'correctly constructed'
9 9 schemes (networks), the latter usually being defined by

fY ly . . t
J 2 m fy
induction.
References
Fig. 5. Fig. 6. [1] SHANNON,C. and MACCARTHY, J. (EDS.): Automata studies,
Princeton Univ. Press, 1956.
[2] GLUSHKOV, V.M.: SyntheSiS of numerical automata, Moscow,
This scheme specifies the structural automaton
1962 (in Russian).
<An, S, B m, </>, 1/;>, where S=SI X ... XS/, while the V.A. Buevich
functions </> and I/; are defined as follows: S. V. Aleshin
<f>SI, ... ,s/), (a I> ... ,an = (S'I, ... ,s~), Editorial comments. See the Editorial Comment to the
article Automata, equivalence of.
Sj = cf>j(S;' j;(a I> , an, 1/11 (S I), ... ,I/I/(s/)));
I/;(SI> ,s/), (a I> ... ,an = (bl> ... ,bm ), AMS 1980 Subject Classification: 68005

b; = gj(al>'" ,an,I/II(SI), ... ,I/I/(s/. AUTOMATA, MINIMIZATION OF - Minimization

An important example of structural automata are
logical networks (cf. Automaton, finite). F~Ee _6 ~h<?ws
the canonical scheme of an automaton <A, S, B, </>, 1/;>
isomorphic to _ th~ au!omaton ~ (cf. Fig. 3 for its
diagram), A =S =B = {O, 1}. The automaton
21" = <s, q;,
S, S, 1/;' > is a Moore automaton such that
q;
</>'(1, z)=z, (0, z)=z.
Structural automata are frequently described by
means of canonical equations, i.e. systems of the form
ult+ I) = ji(UI(t), ... ,u,(t), x\(t), ... ,xn(t,
YI(t) = g\(UI(t), ... ,U,(t),XI(t), ... ,Xn(t,
of the parameter values of automata which results in
equivalent and, in a sense, optimal automata. The
problem of minimization of automata arises during the
synthesis of automata, and its specific nature will
depend on the approach adopted in the study of the
problem.
If the macro-approach is adopted, the usual pro-
cedure is to minimize the number of states of automata,
thus obtaining minimal or reduced automata. The actual
method of search for reduced automata will depend on
the form in which they are defined and on their type of
behaviour. If, for example, the behaviour of a finite
automaton (cf. Automaton, finite), defined with the aid
of, for example, a transition diagram, IS understood to
mean the system of finitely-determined functions (cf.
Finitely-detennined function) realized by the automa-
ton, the minimal finite automaton equivalent to the ini-
tial automaton can be effectively found. This is based
on Moore's theorem, according to which it may be esta-
blished by an experiment of length n - I whether or not
two states of an automaton with n states are distin-

288

guishable from one another. The mmmuzation algo-
rithm consists of constructing an automaton as defined
above whose states correspond to the classes of indis-
tinguishable states of the initial automaton, whereas the
transition and the output functions are defined by
means of representatives of these classes. The minimal
automaton is unique up to a state-isomorphism. If the
final automaton is considered as an acceptor represent-
ing - by means of the subset of isolated states - an
event described with the aid of the given regular
expression, the minimal automaton is effectively con-
structed, and the algorithm of its construction may be
divided into two stages. In its first stage, the initial reg-
ular expression is used to construct some automaton,
not necessarily minimal, representing the corresponding
regular event. Such an automaton may be constructed,
for example, by induction from the operations of union
U, concatenation 0 and iteration * which are included
in the regular expression. The automata m1, m 2 and m3
which represent, respectively, the events R 1, R2 and R3
serve to effect the special constructions of the automata
m4, ms and m6 which represent the events R 1U R 2,
R 1 0R 2 and R;; moreover, the number of states of m
does not exceed the product of the number of states of
m1and m2; the number of states of ms does, generally,
not exceed the product of the number of states of m1
with the number of all subsets of the set of states of
m2; while the number of states of m6 does not exceed
the number of subsets of the set of states of m 3 During
the second stage of the algorithm the number of states
of the resulting automaton is minimized in the usual
way, the classes of indistinguishable final states of the
initial automaton corresponding to the final states of
the minimal automaton. The minimal automaton
representing a given super-event is constructed in the
same manner. There are unique (up to a state-
isomorphism) minimal automata which represent given
events and super-events.
A reduced automaton, which is unique up to a state-
isomorphism, also exists for non-deterministic finite
automata, and for both deterministic and non-
deterministic infinite automata. In the former case, this
automaton is found using a method similar to that in
the case of finite automata discussed above, while in
the second case its construction is, as a rule, non-
effective. For a stochastic automaton with a finite
number of states there exists, generally speaking, a con-
tinuum of various reduced automata are equivalent to
the initial automaton. This is why there is no effective
way of describing the set of reduced automata derived
from a given stochastic automaton.
In the micro-approach to the study of finite auto-
mata, a given automaton is constructed starting from
some finite selection of a basis ~ of logical networks.
AUTOMATA, THEORY OF
Given an automaton m, defined, for example, by
canonical equations, it is required to exhibit a logical
network S, constructed from the elements of the basis
~ using superposition and feedback operations, realiz-
ing mand containing a minimum number L(m) of ele-
ments of ~ necessary for the realization of this automa-
ton. In this sense the network S will be the optimal
scheme of the automaton m. In the general case, the
problem of realizing an arbitrary automaton mwith the
aid of networks over a basis ~ is unsolvable, and for
this reason the complexity function L(m) of the auto-
maton mcannot be calculated. If a basis of mis known
to be complete (cf. Automata, complete systems of), the
construction of any optimal network for m can be
effectively realized e.g. in the following manner. It is
known that the realizability check of an automaton m
with the aid of a given network S can be effectively
established. Moreover, for a given number I the number
of networks over a base ~ of complexity at most I can
be calculated, and all these networks can be effectively
constructed. Therefore, a sequential inspection of all
networks of complexity 1, ... ,L(m) for the realizabil-
4 ity of the automaton m may be used as the algorithm
for constructing all optimal networks for a given auto-
maton m. The problem of the behaviour of the function
L(m) and of certain of its generalizations is part of the
general problem of synthesis of automata. There exist
several heuristic algorithms for the synthesis of minimal
schemes for automata, based on the special properties
of the bases and on specific interpretations of the con-
cept of optimality. For references see Automaton, fm-
ite.
V.A. Buevich
AMS 1980 Subject Classification: 68005
AUTOMATA, THEOR.Y OF - A branch of the theory
of control systems whose subject is the study of
matematical models of transformers of discrete infor-
mation, known as automata. In a sense, such
transformers may be both real mechanisms (computers,
automata, living organisms, etc.) and abstract systems
(mathematical machines, axiomatic theories, etc.). The
theory of automata was born in the mid-twentieth cen-
tury in connection with finite automata (d. Automaton,
fmite), which are mathematical models of nervous sys-
tems and electronic computers. The class of objects
and the scope of problems dealt with by the theory of
automata was subsequently enlarged, and came to
include certain concepts in and problems of other
branches of mathematics. The theory of automata is
most closely connected with the theory of algorithms
(d. Algorithms, theory of), particularly so with the
theory of abstract machines, since automata can be
regarded as a special case of such machines.
289
AUTOMATA, THEORY OF
Most problems of the theory of automata are similar
to those of the other main types of control systems.
They include the problems of analysis and synthesis of
automata, completeness problems, minimization prob-
lems, problems on equivalent transformations of auto-
mata, etc. The problems of analysis concern the
description of the behaviour of a given automaton or
the determination of certain of its properties from its
behaviour and from given incomplete data on its
behaviour (cf. Automaton, behaviour of an; Automata,
experiments with). Problems in the synthesis of auto-
mata include the construction of automata with pre-
defined properties or automata functioning in a desired
manner (cf. Synthesis problems). Related problems con-
cern estimates of the complexity of automata behaving
in a given manner and the construction of automata
which are, in some sense, optimal (cf. Automata, minim-
ization of). In addition, there arises the problem of
completeness in connection with classes of initialized
automata or mappings of automata (cf. Functional sys-
tem; Automata, complete systems of). The problem of
equivalent transformations may be posed both for auto-
mata and in various problems concerning their
behaviour (cf. Equivalent transformations). In addition
to problems of the above type, which are common to
many control systems, there are problems specific to
automata. Thus, depending on the particular problem
at hand, it may be convenient to specify the behaviour
of automata in various languages (regular expressions,
canonical equations, language of logical predicates, etc.,
cf. Automata, methods of specification of), which gives
rise to problems concerning the choice of an adequate
and convenient language, and concerning the
translation from one language to another. Other prob-
lems, closely connected with synthesis and equivalent
transformation problems, is the problem of minimiza-
tion of the number of states of an automaton and that
of obtaining corresponding estimates. Fairly simple
algorithms have been developed for finite automata
making it possible, with the aid of regular expressions,
to obtain automata with the minimum possible number
of states and representing appropriate events (cf. Auto-
mata, minimization of). A related set of problems con-
cerns the modelling of the behaviour of automata in
one class by automata in another class. Problems of the
minimization of modelling automata and estimation of
their complexity are also of interest. For instance, on
passing from a non-deterministic automaton (source)
representing (generating) a regular set of words to a
finite automaton representing the same set, the number
of states may increase exponentially. A special field of
the theory of automata are the so-called experiments
with automata (cf. Automata, experiments with). Here,
the principal problem is to obtain certain information
290
about the structure of an automaton by observing its
reactions to certain outside effects. This involves a wide
circle of problems concerning the classification of
experiments and problem-solving by certain kinds of
experiments, as well as estimates of the shortest experi-
ment sufficient to solve a given problem. The concept
of experiments with automata is also used in control
problems concerning automata (cf. Reliability and
inspection of control systems). Problems of the interac-
tion of automata - either with each other or with
given external environments - are dealt with in
branches dealing with automaton games and the
behaviour of automata (cf. Automaton, behaviour of an).
Many of these problems may be considered as being
algorithmic in nature. Most of them have a positive
solution for finite automata.
The theory of automata has found applications both
in other branches of mathematics and in solving practi-
cal.problems. Thus, it can be used to prove the solva-
bility of certain formal calculi. The application of the
methods and concepts of the theory of automata to the
study of formal and natural languages gave rise to a
special branch of the theory of algorithms: mathemati-
cal linguistics. The concept of an automaton can serve
as model in a great variety of problems, and can thus
be employed in many studies, both in pure and applied
sciences.
References
[1) SHANNON, C. and MACCARTHY, J. (EDS.): Automata studies,
Princeton Univ. Press, 1956.
[2) GLUSHKOV, Y.M.: SyntheSiS of numerical automata, Moscow,
1962 (in Russian).
[3) ThAKHTENBROT, B.A. and BARZDIN', YA.M.: Finite automata.
Behaviour and synthesis, North-Holland, 1973 (translated from
the Russian).
VB. Kudryavtsev
Yu.l. Yanov
Editorial comments. This article refers to only some
aspects of the theory. One should also consult the articles
Formal languages; L-systems; Complexity.
AMS 1980 Subject Classification: 680XX, 93AXX,
03005, 68F10
AUTOMATIC CONTROL THEORY - The science
dealing with methods for the determination of laws for
controlling systems that can be realized by automatic
devices. Historically, such methods were first applied
to proce~~t:~ which ",-ere mainly technical in nature [1].
Thus, for instance, an aircraft in flight is a system, the
control laws of which ensure that it remains on the
required trajectory. The laws are realized by means of a
system of transducers (measuring devices) and actua-
tors, which is known as the automatic pilot. This
development was due to three reasons: many control
systems had been identified by classical science (to
identify a control system means to write down its

mathematical model, e.g. relationships such as (I) and
(2) below); long before the development of automatic
control theory, thanks to the knowledge of the funda-
mental laws of nature, there was a well-developed
mathematical apparatus of differential equations and
especially an appparatus for the theory of steady
motion [2]; engineers had discovered the idea of a feed-
back law (see below) and found means for its realiza-
tion.
The simplest control systems are described by an
ordinary (vector) differential equation
x = j(x, u, t) (1)
and an inequality
N(x, u, t) ;;;. 0, (2)
where x {x I, . . . ,xn } is the state vector of the system,
u {U I, . . . ,u,} is the vector of controls which can be
suitably chosen, and t is time. Equation (1) is the
mathematical representation of the laws governing the
control system, while the inequality (2) establishes its
domain of definition.
Let U be some given class of functions u (t) (e.g.
piecewise-continuous functions) whose numerical values
satisfy (2). Any function U(t)E U will be called a per-
missible control. Equation (1) will be called a
mathematical model of the control system if:
1) A domain N(x, u, t);a.O in which the function
f (x, u, t) is defined has been specified;
2) A time interval .r= [ti' tf] (or [ti' tf)' if tf= 00) dur-
ing which the motion x(t) is observed, has been speci-
fied;
3) A class of permissible controls has been specified;
4) The domain N;;;'O and the function f(x, u, t) are
such that equation (1) has a unique solution defined for
any t E.'T, Xo EN, whatever the permissible control u (t).
Furthermore, f (x, u, t) in (1) is always assumed to be
smooth with respect to all arguments.
Let Xi =x (t;) be an initial and let xf=x (tf) be a
(desired) final state of the control system. The state xf
is known as the target of the control. Automatic control
theory must solve two major problems: the problem of
programming, i.e. of finding those controls u(t) permit-
ting the target to be reached from Xi; and the determi-
nation of the feedback laws (see below). Both prob-
lems are solved under the assumption of complete con-
trollability (1).
The system (1) is called completely controllable if, for
any Xi, XfEN, there is at least one permissible control
u(t) and one interval .rfor which the control target is
attainable. If this condition is not met, one says that
the object is incompletely controllable. This gives rise to
a preliminary problem: Given the mathematical model
(1), find the criteria of controllability. At the time of
AUTOMATIC CONTROL THEORY
wntmg (1977) only insignificant progress has been
made towards its solution. If equation (1) is linear
x = Ax+Bu, (3)
where A, B are stationary matrices, the criterion of
complete controllability is formulated as follows: For
(3) to be completely controllable it is necessary and suf-
ficient that the rank of the matrix
Q = II B AB ... A n -I B II (4)
be n. The matrix (4) is known as the controllability
matrix.
If A, B are known differentiable functions of t, the
controllability matrix is given by
(5)
where
dL -
LI(t) = B(t), Lk(t) = A(t)Lk - I -----;;;-,
k 1
k=2, ... ,no
The following theorem applies to this case: For (3) to
be completely controllable it is sufficient if at at least
one point t* E.rthe rank of the matrix (5) equals n [3].
Criteria of controllability for non-linear systems are
unknown (up to 1977).
The first principal task of automatic control theory is
to select the permissible controls that ensure that the
target xf is attained. There are two methods of solving
this problem. In the first method, the chief designer of
the system (1) arbitrarily determines a certain type of
motion for which the target xf is attainable and selects
a suitable control. This solution of the programming
problem is in fact used in many instances. In the
second method a permissible control minimizing a
given cost of controls is sought. The mathematical for-
mulation of the problem is then as follows. The data
are: the mathematical model of the controlled system
(1) and (2); the boundary conditions for the vector x,
which will be symbolically written as
(i, j) = 0; (6)
a smooth function G(x, t); and the cost of the controls
used
t:J.G
If= G(x, t) ;- (7)
The programming problem is to find, among the permis-
sible controls, a control satisfying conditions (6) and
yielding the minimum value of the functional (7).
Necessary conditions for a minimum for this non-
classical variational problem are given by the L.S. Pon-
tryagin 'maximum principle' [4] (cf. Pontryagin max-
imum principle). An auxiliary vector tPN]' ... , tPn}
and the auxiliary scalar function
H(l/I, x, u, t) = l/I-j(x, u, t) (8)
are introduced. The function H makes it possible to
write equation (1) and an equation for the vector tP in
the following form:
291
AUTOMATIC CONTROL THEORY
. aH' aH
X = ai' I/; = aX (9)
Equation (9) is linear and homogeneous with respect to
I/; and has a unique continuous solution, which is
defined for any initial condition 1/;{ti ) and t EY. The vec-
tor I/; will be called a non-zero vector if at least one of
its components does not vanish for t EY. The following
theorem is true: For the curve xo, Uo to constitute a
strong minimum of the functional (7) it is necessary
that a non-zero continuous vector 1/;, as defined by
equation (9), exists at which the function H(I/;, x, u, t)
has a (pointwise) maximum with respect to u, and that
the transversality condition
[8G-H8t+~l/;a8xay = 0
is met. Let XO(t, Xi' Xj), Uo(t, Xi, Xj) be solutions of the
corresponding problem. It has then been shown that
for stationary systems the function H (1/;0, xo, UO) satis-
fies the condition
H(I/;,xO,UO) = C, (10)
where C is a constant, so that (10) is a first integral.
The solution uO, XO is known as a program control.
Let uO, XO be a (not necessarily optimal) program
control. It was found that the knowledge of only one
program control is insufficient to attain the target. This
is because the program uO, XO is usually unstable with
respect to, however small, changes in the problem, in
particular to the most important changes, those in the
initial and final values (i, f) or, in other words, the
problem is ill-posed. However, this ill-posedness is such
that it can be corrected by means of automatic stabili-
zation, based solely on the 'feedback principle'. Hence
the second main task of control: the determination of
feedback laws.
Let y be the vector of disturbed motion of the system
and let ~ be the vector describing the additional deflec-
tion of the control device intended to quench the dis-
turbed motion. To realize the deflection ~ a suitable
control source must be provided for in advance. The
disturbed motion is described by the equation:
y = Ay+B~+cJ>(y,~,t)+f'(t). (II)
where A and B are known matrices determined by the
motion of xv, uo, and are assumed to be known func-
tions of the time: cf> i." the non-linear part of the
development of the function f(x, u, t); f'(t) is the con-
stantly acting force of perturbation, which originates
either from an inaccurate determination of the pro-
grammed motion or from additional forces which were
neglected in constructing the model (l)._Equation ill)
is defined in a neighbourhood II f II ~H,. where H is
usually quite smalL but in certain cases may be any fin-
ite positive number or even 00.
It should be noted that, in general, the fact that the
292
system (1) is completely controllable does not mean
that the system (11) is completely controllable as well.
One says that (11) is observable along the coordinates
flo ... ,f" r~n, if one has at his disposal a set of
measuring instruments that continuously determines the
coordinates at any moment of time t EY. The signifi-
cance of this definition can be illustrated by consider-
ing the longitudinal motion of an aircraft. Even though
aviation is more than 50 years old, an instrument that
would measure the disturbance of the attack angle of
the aircraft wing or the altitude of its flight near the
ground has not yet been invented. The totality of meas-
ured coordinates is called the field of regulation and is
denoted by P(y], ... ,Yr), r~n.
Consider the totality of permissible controls ~, deter-
mined over the field P:
~ = ~], ... ,f" t,p), r~n, (12)
where p is a vector or a matrix parameter. One says
that the control (12) represents a feedback law if the
closure operation (i.e. substitution of (12) into (11
yields the system
y= Ay+B~(y], ... ,y"t,p)+ (l3)
+cJ>(y'~(yI"" ,y,,!,p),!),
such that its undisturbed motion y =0 is asymptotically
stable (cf. Asymptotically-stable solution). The system
(13) is said to be asymptotically stable if its undisturbed
motion y = 0 is asymptotically stable.
There are two classes of problems which may be for-
mulated in the context of the closed system (13): The
class of analytic and of synthetic problems.
Consider a permissible control (12) given up to the
selection of the parameter p, e.g.:
PII Plr
~=
prl prr
Yr
The analytic problem: To determine the domain S of
values of the parameter p for which the closed system
(13) is asymptotically stable. This domain is con-
structed by methods developed in the theory of stability
of motion (cf. Stability theory), which is extensively
employed in the theory of automatic control. In partic-
ular, one may mention the methods of frequency
analysis; methods based on the first approximation to
Lyapunov stability theory (the theorems of Hurwitz,
Routh, etc.), on the direct Lyapunov method of con-
structing v-functions, on the Lyapunov- Poincare
theory of constructing periodic solutions, on the
method of harmonic balance, on the B.V. Bulgakov
method, on the A.A. Andronov method, and on the
theory of point transformations of surfaces [5). The last
group of methods makes it possible not only to con-

struct domains S in the space P, but also to analyze the
parameters of the stable periodic solutions of equation
(13) which describe the auto-oscillatory motion of the
system (13). All these methods are widely employed in
the practice of automatic control, and are studied in
the framework of various specifications in schools of
higher learning [5].
If S is non-empty, a control (12) is called a feedback
law or regulation law. Its realization, which is effected
using a system of measuring instruments, amplifiers,
converters and executing mechanisms, is known as a
regulator.
Another problem of considerable practical
importance, which is closely connected with the ana-
lytic problem, is how to construct the boundary of the
domain of attraction [6], [7]. Consider the system (13) in
which pES. The set of values y(ti)=yo containing the
point y =0 for which the closed system (13) retains the
property of asymptotic stability, is known as the
domain of attraction of the trivial solution y =0. The
problem is to determine the boundary of the domain of
attraction for a given closed system (13) and a point
pES.
Modem scientific literature does not contain effective
methods of constructing the boundary of the domain of
attraction, except in rare cases in which it is possible to
construct unstable periodic solutions of the closed sys-
tem. However, there are certain methods which allow
one to construct the boundary of a set of values of yo
totally contained in the domain Qf attraction. These
methods are based in most cases on the evaluation of a
domain in phase space in which the Lyapunov function
satisfies the condition v ~O, V';;;;O [7].
Any solution y(t,yo,p) of the closed system (13)
represents a so-called transition process. In most cases
of practical importance the mere solution of the stabil-
ity problem is not enough. The development of a pro-
ject involves supplementary conditions of considerable
practical importance, which require the transition pro-
cess to have certain additional features. The nature of
these requirements and the list of these features are
closely connected with the physical nature of the con-
trolled object. In analytic problems it may often be
possible, by a suitable choice of the parameter p, to
preserve the desired properties of the transition process,
e.g. the pre-set regulation time t. The problem of
choosing the parameter p is known as the problem of
quality of regulation [5], and methods for solving this
problem are connected with some construction of esti-
mates for solutions y(t,yo,p): either by actual integra-
tion of equation (13) or by experimental evaluation of
such solutions with the aid of an analogue or digital
computer.
The analytic problems of transition processes have
AUTOMATIC CONTROL THEORY
many other formulations in all cases in which r(t) is a
random function - in servomechanisms for example,
[5], [8]. Other formulations are concerned with the pos-
sibility of a random alteration of the matrices A and B
or even of the function cp [5], [8]. This gave rise to the
development of methods for studying random
processes, methods of adaptation and learning
machines [9]. Transition processes in systems with delay
mechanisms and with distributed parameters (see [10],
[11]) and with a variable structure (see [12]) have also
been studied.
The synthesis problem: Given equation (11), a field of
regulation P(y], ... ,y,), r';;;;n, and a set
~], ... ,y" t) of permissible controls, to find the
whole set M of feedback laws [13]. One of the most
important variant of this problem is the problem of the
structure of minimal fields. A field P(y], ... ,r,), r ';;;;n,
is called a minimal field if it contains at least one feed-
back law and if the dimension r of the field is minimal.
The problem is to determine the structure
P(Ya" ... ,YaJ of all minimal fields for a given equa-
tion (11) and a set of permissible controls. The
following example illustrates the nature of the problem:
z = Az+Bu,
o mOO 00
-m 0 0 n 1 0
A o OOk,B= o0
o -n k 0 o1
where m, n, k are given numbers. The permissible con-
trols are the set of piecewise-continuous functions
U2, U4 that take their values from the domain
Iu21 .;;; U2, Iu41 .;;; U4-
The minimal fields in this problem are either the field
P(Z2) or the field P(Z4). The dimension of each field is
one and cannot be reduced [13].
So far (1977) only one method is known for the syn-
thesis of feedback laws; it is based on Lyapunov func-
tions [13]. A relevant theorem is that of
Barbashin - Krasovskil [6], [10], [15]: In order for the
undisturbed motion y = 0 of the closed system
j = cp(y) (14)
to be asymptotically stable, it is sufficient for a
positive-definite function v (y) to exist, such that by
equation (14) its complete derivative is a function w(y)
which is semi-definite negative, and such that on the
manifold w (y) = 0 no complete trajectory of the system
(14), except for y =0, lies. The problem of finding out
about the existence and the structure of the minimal
fields is of major practical importance, since these
fields determine the possible requirements of the chief
designer concerning the minimum weight, complexity
293
AUTOMATIC CONTROL THEORY
and cost price of the control system and its maximum
reliability. The problem is also of scientific and practi-
cal interest in connection with infinite-dimensional sys-
tems as encountered in technology, biology, medicine,
economics and sociology.
In designing control systems it is unfortunately
impracticable to restrict the work to solving problems
of synthesis of feedback laws. In most cases the
requirements of the chief designer are aimed at securing
certain important specific properties of the transition
process in the closed system. The importance of such
requirements is demonstrated by the importance of
monitoring an atomic reactor. If the transition process
takes more than 5 seconds or if the maximum value of
some of its coordinates exceeds a certain value, an
atomic explosion follows. This gives rise to new prob-
lems of synthesis of regulation laws, based on the set
M. Below the formulation of one such problem is given.
Consider two spheres Ilyo I =R, YO=Yi, Ily(td II =f;
Rf are given numerical values. Now consider the
set M of all feedback laws. The closure by means of
any of them yields the equation:
j = Ay+BYh'" ,y" t)+
(15)
</>{y, Y l, ... ,y" t), t).
Consider the entire set of solutions y(t, Yo) of equation
(15) which start on the sphere R and call them R-
solutions. Since the system is asymptotically stable for
any Yo on the sphere, there exists a moment of time t1
during which the conditions
Ily(tl,yo) I = (, Ily(t,yo) II < (,
are valid for any t>t1.
Let
t* = supt l .
Yo
It is clear from the definition of t 1 that t * exists. The
interval t * - ti is called the regulation time (the time of
damping of the transition process) in the closed system
(15) if any arbitrary R-solution starts at the f-sphere if
t1 "';;;t*, but remains inside it if t>t1. It is clear that the
regulation time is a functional of the form
t * =t * (R, f, ~). Let T be a given number. There arises
the problem of synthesis of fast-acting regulators: Given
a set M of feedback laws, one has to isolate its subset
M 1 on which the regulation time in a closed system
satisfies the condition
t* -t, ~ T
One can formulate in a similar manner the problems of
synthesis of the sets M 2, . . . ,Mk of feedback laws,
which satisfy the other k -1 requirements of the chief
designer.
The principal synthesis problem of satisfying all the
requirements of the chief designer is solvable if the sets
M 1 , ,Mk have a non-empty intersection [13].
294
The synthesis problem has been solved in greatest
detail for the case in which the field P has maximal
dimension, r = n, while the cost index of the system is
characterized by the functional
J = fw(y,~,t)dt, (16)
where w(y,~, t) is a positive-definite function of y, f
The problem is then known as the problem of analytic
construction of optimum regulators [14] and is in fact
thus formulated. The data include equation (11), a class
of permissible controls g(y, t) defined over the field
P (y) of maximal dimension, and the functional (16).
One is required to find a control g=g(y, t) for which
the functional (16) assumes its minimum value. This
problem is solved by the following theorem: If equation
(11) is such that it is possible to find an upper semi-
continuous positive-definite function v(y, t) and a
function go(y, t) such that the equality
avO avO
--at+ ay(Ay+B~+</>{y,~,t=O (17)
is true, and the inequality
avO avO
--at + ay (Ay + B~+</>{y,~, t ;;;. O.
00
V(to,yo) = min fw(y,~, t)dt,
~
is true. The function v(y, t) is known as the Lyapunov
optimum function [15]. It is a solution of the partial dif-
ferential equation (17), of Hamilton - Jacobi type, satis-
fying the condition v (y( (0), (0) = O. Methods for the
effective solution of such a problem have been
developed for the case in which the functions wand cp
can be expanded in a convergent power series in y, g
with coefficients which are bounded continuous func-
tions of t. Of fundamental importance is the solvability
of the problem of linear approximation to equation (11)
and the optimization with respect to the integral of
only second-order terms contained in the development
of w. This problem is solvable if the condition of com-
plete controllability is satisfied [15].
References
[I] MAKSVELL, D.K., VISHNEGRADSKli. I.A. and STODOLA. A.: The
theory of automatic regulation, Moscow, 1942 (in Russian).
[2] CHETAEV, N.G.: Stabili~v of motion, Moscow, 1965 (in Rus-
sian).
[3] KRASOVSKII. N.N.: Theory of control of motion. Moscow. 1968
(ill Russian).
[4] PONTRYAGIN. L.S., ET At.: The mathematical theorv of optimal
processes. Moscow. 1969 (in Russian).
[5] Technical cybernetics. Moscow, 1967.
[6] BARBASHIN, E.A.: inrroduction to the theon' of stabilitr.
Wolters-Noordhoff, 1970 (translated from the Russian).
[7] ZUBOV. V.1.: Mathematical methods of stu ding svstem.1 of
automatic control. Leningrad. 1959 (in Russian).
[R] PuGACHEV, V.S.: The theory' random functions and its applica-
tions to problems of automatic control, Moscow, 1962 (in Rus-
sian).

[9] TSYPKIN, Y A.Z.: Foundations of the theory of serving systems,
Moscow, 1970 (in Russian).
[10] KRAsovSKli, N.N.: Stability of motion. Applications of
Lyapunov's second method to difJerential systems and equations
with delay, Stanford Univ. Press, 1963 (translated from the
Russian).
[II] BESEKERSKIi, V.A. and PoPov, E.P.: The theory of automatic
control systems, Moscow, 1966 (in Russian).
[12] EMELYANOV, S.V. (ED.): Theory of systems with a variable struc-
ture, Moscow, 1970 (in Russian).
[13] LETOV, A.M.: 'Some unsolved problems in control theory', Dif-
ferential equations N. Y 6, no. 4 (1970), 455-472.
(DifJerentsial'nye Uravneniya 6, no. 4 (1970),592-615)
[14] LETOV, A.M.: Dynamics offlight and control, Moscow, 1969 (in
Russian).
[15] MALKIN, 1.0.: Theorie der Stabilitiit einer Bewegung, R. Olden-
burg, Miinchen, 1959 (translated from the Russian).
A.M. Letov
Editorial comments. The article above reflects a different
tradition and terminology than customary in the non-Russian
literature. It also almost totally ignores the vast amount of
important results concerning automatic (optimal) control
theory that have appeared in the non-Russian literature.
Rougly speaking a control system is a device whose
future development (either in a deterministic or stochastic
sense) is determined entirely by its present state and the
present and future values of the control parameters. The
'recipe' which determines future behaviour can be almost
anything, e.g. a differential equation
x = t(x, u, t), xER n , uER m, tER, (A1)
as in the main article above, or a difference equation
Xt +1 = teXt, Ut, t), XtER n , utER m , t=O, 1, ... , (A2)
or more generally, a differential equation with delays, an
evolution equation in some function space, a more general
partial differential equation, ... , or any combination of
these. Frequently there are constraints on the values the
control uER m may take, e.g., I U II ~M where M is a con-
stant. However there certainly are (engineering) situations
when the controls u(t) which can be used at time t depend
both explicitly on time and the current state xU) of the sys-
tem (which of course in turn depends on the past controls
used). This may make it complicated to describe the space
of admissible or permissible controls in simple terms.
It may even be not very well possible to describe the con-
trol structure in such terms as equations (A 1) and con-
straints N(x, u, t);;:.0 defined on Rn X RmX R. This happens,
e.g., in the case of the control of an artificial satellite where
the space Rm in which the controls take their value may
depend on the state x, being, e.g., the tangent space to a
sphere at the point x. Then controls become sections in a
vector bundle usually subject to further size constraints.
At this point some of the natural and traditional problems
concern reachability, and control/ability and state space
feedback. Give an initial state x(O) the control system is said
to be completely reachable (from x(O if for all states x
there is an admissible control steering x(O) to x. This is
termed complete controllability in the article above. Control-
lability in the Western literature is usually reserved for the
opposite notion: Given a (desired) final state x(f) the sys-
AUTOMATIC CONTROL THEORY
tem is said to be completely control/able (to x(f if for
every possible initial state x there is a control which steers x
to x(f). In the case of continuous-time, time-invariant,
finite-dimensional systems the two notions coincide, but this
is not always the case.
There are also many result on reachability and controlla-
bility for non-linear systems, especially for local controllabil-
ity and reachability. These results often are stated in terms
of Lie algebras associated to the control system. E.g., for a
control system of the form
. m
X = f(x) + ~u;g;(x), xER n , uER m, (A3)
;=1
one considers the Lie algebra .P(~) spanned by the vector
fields (ad f)k(g;) , i=1, ... ,m; k=0,1, ... , where
ad f(g)=[f, g], (ad f)k=(ad f)o(ad f)k-1. Given a Lie algebra
of vector fields .P on a manifold M the rank of .P at a point
xEM, rkx.P' is the dimension of the space of tangent vec-
tors {X(x): XE.P} C TxM. A local reachability result now
says that rkx .P(~) = n implies local reachability at x. There
are also various global and necessary-condition-type results
especially in the case of analytic systems. Cf. [A2], [A 10],
[A14] for a first impression of the available results.
Given a control system (A 1) an important general ques-
tion concerns to what extent it can be changed by feed-
back, in this case state space feedback. This means the fol-
lowing. A state space feedback law is a suitable mapping
k: Rn~Rm, xf->u=k(x). Inserting this in (A1) results in the
closed loop system x= f(x, k(x), t). Cf. Fig. 1.
V 1"""'1.--"' U
L
If
Fig. 1.
Often 'new' controls vER m are also introduced and one
considers, e.g. the new control system x = f(x, k(x) + v, t)
or, more generally, x= f(x, k(x, v), t). Dynamic feedback
laws in which the function k(x, v) is replaced by a complete
input-output system (cf. below)
y = g(y, x, v, t), U = hey), (A4)
are also frequently considered. Important questions now
are, e.g., whether for a given system (A1) a k(x) or k(x, v)
or system (A4) of various kinds can be found such that the
resulting closed loop system is stable. Or, whether by
means of feedback a system can be linearized or imbedded
in a linear system (4). Cf. [A6] for a selection of results.
Using (large) controls in engineering situations can be
expensive. This leads to the idea of a control system with
cost functional J which is often given in terms such as
l
t,
J(u) = g(x, u, t) dt+ F(XU1 .
There result optimal control questions such as finding that
admissible function U which minimizes J(u) (and steers x to
295
AUTOMATIC CONTROL THEORY
the target set) (optimal open loop control) and finding a
minimizing feedback control law u=k(x) (optimal closed
loop control). In practice the case of a linear system (4) with
a quadratic criterion is very important, the so called LO
problem. Then
G(x, u, t) = u TRu+2u TSx+X TOx, F(x) = xTMx. (A5)
Here the upper T denotes transpose and R, S, 0, M are suit-
able matrices (which may depend on time as may A and 8).
In this case under suitable positive-definitness assumptions
on R, M and the block matrix
II~ ~II'
the optimal solution exists. It is of feedback type and is
obtained as follows. Consider the matrix Riccati equation
k = -AK-KA+(S+BTKlR-'(S+BTK)-O, (A6)
K(t,) = M.
Solve it backwards up to time to. Then the optimal control is
given by
(A7)
The solution (method) extends to the case where the linear
system is in addition subject to Gaussian stochastic distur-
bances, cf. [A9] , [A 17].
Also for non-linear systems there are synthesis results for
optimal (feedback) control. One can, for instance, use the
Pontyagin maximum principle to determine the optimal
(open loop) control for each initial state xERn (if it exists
and is unique). This yields a mapping x .... u which is a candi-
date for an optimal feedback control law and under suitable
regularity assumptions this is indeed the case, cf. [A3] ,
[A 15]. Some standard treatises on optimal control in various
settings are [A5] , [A 12], [A 13].
In many situations it cannot be assumed that the state x
of a control system (A1) is directly accesible for control pur-
poses, e.g. for the implementation of a feedback law. More
generally only certain derived quantities are immediately
observable. (Think for example of the various measuring
devices in an aircraft as compared to the complete state
description of the aircraft.) This leads to the idea of an
imput-output (dynamicaO system, or briefly (dynamical) sys-
tem, also called plant (cf. Fig. 2),
x= f(x, u, t), y = hex, u, t), (AB)
xERn, UER m, YERP.
u >
Fig. 2.
Here the U E Rm are viewed as controls or inputs and the
YERP as observations or outputs. Let x(t, u; xo) denote the
solution of the first equation of (AB) for the initial condition
x(to)=xa and a given u(t). The system (AB) is called com-
296
pletely observable if for any two initial states xc, X~ ERn and
(known) control u there holds that
y(x(t, u;xo), u, t)=y(x(t, u;x~), u, t) for all t';;il=ta implies
Xo =x~. This is more or less what is meant by the phrase
'observable along the coordinates y" ... ,Yn' in the main
article above. In the case of a time-invariant linear system
x= Ax+Bu, y = Cx, (A9)
complete observability holds only if and only the (block)
observabi/ity matrix" CTCTA T ... CT(AT)n-,,, is of full
rank n. This is completely dual to the reachability (controla-
bility) result for linear systems in the article above.
In this setting of input-ouput systems many of the prob-
lems indicated above acquire an output analogue, e.g. out-
put feedback stabilization, where (in the Simplest case) a
x
function u = key) is sought such that = f(x, k(h(x, u, t), t
is (asymptotically) stable, the dynamic output feedback
problem and the optimal output feedback control problem.
In addition new natural questions arise such as whether it is
possible by means of some kind of feedback to make cer-
tain outputs independent of certain inputs (decoupling prob-
lems). In case that there are additional stochastic distur-
bances possibly both in the evolution of x and in the meas-
urements y problems of filtering are added to all this. E.g.
the problem of finding the best estimate x(t) of the state of
the system given the observations yes), ta o;;;;so;;;; t, [A9].
There is also the so-called realization problem. Given an
initial state xa a system such as (AB) defines a mapping
from a space of input functions u to a space of output func-
tions y, and the question arises which mappings can be
realized by means of systems such as (AB). Cf. [A6] for a
survey of results in this direction in the non-linear deter-
ministic case and [A9] for results in the linear stochastic
case.
With the notable exception of the output feedback prob-
lem one can say at the present time that the theory of linear
time-invariant finite-dimensional systems, possibly with
Gaussian noise and quadratic cost criteria, is in a higly-
satisfactory state. Cf., e.g., the standard treatise [A 11]. Gen-
eralization of this substantial body of results to a more gen-
eral setting seems to require sophisticated mathematics, e.g.
algebraic geometry and algebraic K-theory in the case of
families of linear systems [AB] , [A 16], functional analysis
and contraction semi-groups for infinite-dimensional linear
systems [A4], functional analysis, interpolation theory and
Fourier analysis for filtering and prediction [A9] , and folia-
tions, vector bundles, Lie algebras of vector fields and other
notions from differential topology and geometry for non-
lineai systems theory [A2], [A9].
A great deal of research at the moment is concerned with
systems with unknown (or uncertain) parameters. Here
adaptive control is important. This means that one attempts
to design e.g. output feedback control laws which automati-
cally adjust themselves to the unknown parameters.
A good idea of the current state of the art of system and
control theory can be obtained by studying the proceedings
of the yearly IEEE CDC (Institute of Electronic and Electical
Engineers Conference on Decision and Control) confer-

ences and the biyearly MTNS (Mathematical Theory of Net-
works and Systems) conferences.
References
[A 1] BARNETI, S.: Introduction to mathematical control theory,
Oxford Univ. Press, 1975.
[A2] BROCKETI, R.W.: 'Nonlinear systems and differential
geometry', Proc. IEEE 64 (1976), 61-72.
[A3] BRUNOVSKY, P.: 'On the structure of optimal feedback sys-
tems', in Proc. Internat. Congress Mathematicians Helsinki,
1978, Vol. 2, Acad. Sci. Fennica, 1980, pp. 841-896.
[A4] CURTAIN, R.F. and PRITCHARD, AJ.: Infinite-dimensional
linear system theory, Springer, 1978.
[A5] FLEMING, W.H. and RISHEL, R.W.: Deterministic and sto-
chastic optimal control, Springer, 1975.
[A6] FLIEss, M. and HAzEWINKEL, M. (EDS.): Algebraic and
geometric methods in nonlinear control theory, Reidel, 1986.
[A7] HAZEWINKEL, M.: 'On mathematical control engineering',
Gazette des Math. 28, July (1985), 133-151.
[A8] HAZEWINKEL, M.: '(Fine) moduli spaces for linear systems:
What are they and what are they good for', in C.1. Byrnes
and C.F. Martin (eds.): Geometric methods for linear system
theory, Reidel, 1980, pp. 125-193.
[A9] HAZEWINKEL, M. and WILLEMS, J.e. (EDS.): Stochastic
systems: The mathematics of filtering and identification,
Reidel, 1981.
[A 10] JURDJEVIC, V. and KUPKA, I.: 'Control systems on semi-
simple Lie groups and their homogeneous spaces', Ann. Inst.
Fourier 31 (1981), 151-179.
[A11] KWAKERNAAK, H. and SIVAN, R.: Linear optimal control sys-
tems, Wiley, 1972.
[A12] LEE, E.B. and MARKUS, L.: Foundations of optimal control
theory, Wiley, 1967.
[A 13] LIONS, lL.: Optimal control of systems governed by partial
differential equations, Springer, 1971.
[A14] LOHRY, e.: 'Controlabilite des systemes non-lineaires', in
Outils et modeles mathematiques pour I'automatique,
I'analyse des systemes et Ie traitement du signal, Vol. 1,
CNRS, 1981, pp. 187-214.
[A15] SUSSMANN, HJ.: 'Analytic stratifications and control theory',
in Proc. Internal. Congress Mathematicians Helsinki, 1978,
Vol. 2, Acad. Sci. Fennica, 1980, pp. 865-871.
[A16] TANNENBAUM, A.: In variance and system theory: algebraic
and geometric aspects, Springer, 1981.
[A 17] WILLEMS, J.e.: 'Recursive filtering', Statistica Neerlandica 32
(1978), 1-39.
[A18] Handbuch der Systemtheorie, Akademie Verlag, 1986.
AMS 1980 Subject Classification: 93-XX, 49EXX,
49BXX
AUTOMATIC PROGRAMMING - The use of com-
puters to produce computer programs (cf. Program)
automatically from some initial description which is
closer to the original formulation of the problem. The
meaning of 'automatic programming' has shifted in
time, reflecting the general evolution of human-
computer communication and of programming
methods. Originally, automatic programming consisted
in leaving it to the cmputer to produce a translation
into machine code of the description of an algorithm
for solving some problem, expressed in some algo-
rithmic language. Subsequently, automatic programming
was understood to include a method for the directed
synthesis of an algorithm from a compact problem
description, restricted to some fixed class of problems.
AUTOMATIC TRANSLATION
At present (the nineteen-seventies) the concept of
automatic programming is being extended to cover a
comprehensive procedure for the synthesis of an algo-
rithm in the framework of a much wider language for
formulating problems, not limited to a specific class.
Such a synthesis entails a correctness proof for the
algorithm and employs formal approaches such as vari-
ous models for the semantics of natural languages,
predicate calculus and theorem-proving techniques.
A.P. Ershov
Editorial comments.
References
[A1] BALZER, R.: 'A 15-year perspective on automatic program-
ming', IEEE Trans. Software Engineering 11, no. 11 (1985),
1257-1267. Special issue on artificial intelligence and
software engineering, Guest Ed.: J. Mostow.
AMS 1980 Subject Classification: 68RXX
AUTOMATIC TRANSLATION, machine translation -
Translation of texts from one natural language to
another using automatic devices.
Automatic translation is a problem in the area of
modelling and automation of human mental activities
- in the present case, of language. It is performed by
an algorithm, the rules of which do not interact with
the knowledge and intuition of man, i.e. are strictly for-
mal. An algorithm for automatic translation can be
implemented for any suitable automatic device, e.g. by
a digital computer.
Automatic translation is closely connected with the
development of modern structural and mathematical
linguistics, including difficult and fundamental linguis-
tic problems, many of which were neglected in the past
or were not even explicitly stated.
The theoretical basis of the development of
automatic translation is the theory of formal grammars
(cf. Grammar, fonnal). An automatic translation algo-
rithm realizes a certain correspondence between two
languages which are defined by their grammars - a
so-called translation correspondence. It is formulated in
terms of structural descriptions of the properties of
these languages within the framework of their gram-
mars. The most frequently used structural descriptions
are constituent structures and dependency trees (cf.
Syntactic structure). An automatic translation algorithm
accordingly consists of three main parts: 1) a parsing of
the text in the source language, i.e. an analysis of the
structure of the input text, based on the given grammar
of the input language; 2) a transfer, i.e. a transforma-
tion from the structure of the text in the source
language to the structure of the text in the target
language, based on the given translation correspon-
dence; and 3) a synthesis of the text in the target
language, i.e. a transition from the structure of the out-
put text to the specific word sequence.
297
AUTOMATIC TRANSLATION
During the early period of studying automatic trans-
lation the transition rules were usually formulated
directly, in the form of algorithmic rules, without a
preliminary construction of the formal grammars of the
source and target languages, and without an explicit
formulation of the translation correspondence. As a
result, purely linguistic information became confused
with the mathematical formulation of the algorithm. In
a later period, automatic translation algorithms were
constructed on the base of general schemes applicable
to a large number of languages, with the maximum
possible particularization of the translation stages; the
synthetic and the analytical stages were made mutually
independent, with an itemization of all possible
varieties of interpretation of the text, permissible gram-
mars and translation correspondences (so-called multi-
variant translation). Information about the specific
language was in most cases rigidly separated from the
algorithmic part. In this way it became possible to
disregard a large number of details specific for each
language, and to concentrate on the development of
general procedures for obtaining solutions valid under
the specified conditions.
As regards the purely linguistic aspects of automatic
translation, all morphological and almost all syntactic
problems in a given sentence may be considered as
solved. The principal difficulties in the creation of com-
pletely automatic systems of high-quality translation
are due to the relatively backward state-of-the-art of
the semantic theory of languages, which could be used
for an exact formulation of the rules of processing the
sense and the various meanings, and also for under-
standing the logical connection between the separate
sentences of a coherent text. Unless due allowance is
made for the semantics of the text, the translation often
proves ambiguous or faulty. A faulty translation is usu-
ally due to the fact that the translation rules do not
include all possible meanings of the text to be
translated or do not invariably ensure a proper selec-
tion of the meaning, which may depend on a fairly
wide context or even on the knowledge of the subject
matter of the text being translated, on general
knowledge, etc.
Operational industrial systems of high-quality
automatic translations are as yet (nineteen-seventies)
non-existent. There exist experimental systems of vary-
ing degrees of complexity. Simplified and specialized
systems of automatic translation are successfully used
in practical work; these are based on automatic pro-
cessing of scientific and technical information (word-
for-word translation with partial grammatical process-
ing, partial translation and automatic referencing for
on-the-spot information, patent documentation, and for
information search systems). Automatic translation can
29S
serve in the foreseeable future for translating scientific
and technical publications only; automatic translation
of literature and fiction is both unrealistic and unneces-
sary.
References
[1] PANOV, D.Yu.: Automatic translation, Moscow, 1958 (in Rus-
sian). S. vLa. rltW
D" IOV
Editorial comments. Machine translation started in the
USA around 1955.
From a linguistic point of view the programs were naive
(in principle they gave word-by-word translations), and they
were designed in an unstructured way (e.g. linguistic infor-
mation was not separated from the translation algorithm).
The programs were not satisfactory: they were slower, less
accurate and more costly than translations provided by
human translators. Therefore it was decided in the USA in
1966 to stop all support by the government for machine
translation projects. In the years 1966 - 1975 hardly any
fundamental research concerning machine translation was
done. The programs from the first period were developed
further in commercial surroundings, and it are mainly these
programs that are used practically nowadays. Their use is
not in the quality of the translations, but in the saving of
time they give to professional translators when used as auxi-
liary tool (e.g. as text editor). There is up till now only one
fully automatic translation system: the Canadian system
TAUM that translates wheather forecasts from English into
French. Since 1975 there is a revival of research in
machine translation, espeCially in western Europe and
Japan. Important factors in this research are developments
in theoretical linguistics, hardware developments which
change the factors of costs and speed, and the increased
need for specialized translations (especially for technical
manuals).
References
[A1] HUTCHINS. W.J.: Machine translation: past, present, future,
Ellis Norwood limited, Chichister, England, 1986.
[A2] SLOCUM, J.: 'A survey of machine translation: its history,
current status, and future prospects', Computational Linguis-
tics 11 (1985), 1-17.
AMS 1980 Subject Classification: 68F30, 68F05
AUTOMATON - A control system, which is usually
understood to mean a finite automaton (cf. Automaton,
finite) or one of its modifications obtained by changing
components or the mode of operation. The principal
concept - a finite automaton - originated in the
mid-twentieth century in connection wiin aiiempls iu
describe, in mathematical language, the functioning of
nervous systems, of universal computers and of other
real automata. The first such efforts are due to W.
McCulloch and W. Pitts (1943), S.c. Kleene (1951),
A.W. Burks and 1. Wright (1954), and others. A
characteristic feature of such a description is the
discrete character of the corresponding mathematical
models and the finiteness of the domain of their param-

eter values, which accounts for the name 'finite' auto-
maton. The external effects, reactions and states are
considered to be letters of three alphabets named,
respectively, the input alphabet, the output alphabet
and the alphabet of states. The relations governing
their interactions may then be given by two functions
- a transition function and an output function -
which map the pair 'state-input letter' into 'state letter'
and the pair 'state-output letter' into 'output letter',
respectively. At each discrete moment of time the sys-
tem, being in a given state, receives an input signal (a
letter of the input alphabet), emits an output signal (a
letter of the output alphabet, determined by the output
function) and passes into a new state determined by the
transition function. Studies of finite automata are sup-
plemented by studies of their various generalizations
and modifications which reflect various features of real
systems. In the case of a finite automaton
(A, S, B, cp, t/J) the existing modifications may be subdi-
vided into the following three main groups. The first
group includes automata some alphabets A, S or B of
which are infinite, and such automata are said to be
infinite. The second group includes automata in which
the functions t/J and cp are replaced by arbitrary rela-
tionships or by random functions. These are partial,
non-deterministic, probabilistic and other automata.
The third group includes automata with specific sets of
input objects. These are automata with a variable struc-
ture, automata over terms (or tree automata, cf. Auto-
mata, algebraic theory of). There are automata belong-
ing to several groups at the same time, e.g. fuzzy auto-
mata belong to all three groups. Special subclasses of
finite automata are very important; these include, for
instance, memoryless automata, autonomous automata,
reversible automata, etc. In addition, the use of con-
cepts and methods of other branches of mathematics
also generates specific classes of automata and related
problems. For instance, if algebraic methods are
employed there result the concepts of automata over
terms, as well as linear, group, free and other automata
(cf. Automata, algebraic theory of); problems in coding
theory generate the concepts of self-regulating atomata,
reversible automata, etc. (cf. Automaton, probabilistic).
Structural automata also have several generalizations,
mainly consisting in permitting infinite networks and
altering the interconnections between the elements dur-
ing the operation. This leads to the concept of a grow-
ing automaton. The principal modifications and sub-
classes of finite automata, together with their most
important properties, are described below.
Macro approach. l. Infinite automata (first group)
differ from finite automata only by the fact that their
alphabets A, B or S (input, output and set of states)
may be infinite. Most concepts and problems connected
AUTOMATON
with finite automata are naturally extended to infinite
automata. Alphabets of higher cardinality extend the
computational capacities of automata. For instance,
while finite automata realize finitely-determined
(sequential) functions (cf. Finitely-detennined function),
infinite automata may be used to realize any deter-
mined function. Moreover, using infinite automata it is
possible to describe the functioning of any automaton
and machine. At the same time, this very general
nature of infinite automata impairs their significance,
so that most studies concern only special subclasses of
infinite automata, connected with specific models of
control systems.
2. Non-deterministic and asynchronous automata
(second group). A non-deterministic automaton is for-
mally defined as a system (A, S, B, X), where A, S and
B are alphabets in the sense explained above, while
XeS XA X S X B is a transition-to-output relation. If
the relation X is a function mapping S XA into S X B,
the non-deterministic automaton is called a determinis-
tic automaton and is in fact identical with a finite auto-
maton, since in such a case X may be considered as the
pair of functions cp, t/J which map S XA into S and B,
respectively. As distinct from a finite automaton, an
initialized non-deterministic automaton ~S, has several
initial states, which form a subset S I of the set S. The
behaviour of ~S, is usually understood to mean one of
the following generalizations of the behaviour of a fin-
ite automaton.
a) Instead of a function f, an automaton ~S, realizes
a relation / , consisting of all pairs of words
(al ... a", b l . . . b,,)EA * XB* such that there are
states s 1> ,S" + 1> for which S I ES I and
(Sjaisi+lbi)EX is true for any i=l, ... ,no The class of
relations realized by initialized non-deterministic auto-
mata coincides with the class of finitely-determined
relations (cf. Finitely-determined function).
b) An initialized non-deterministic automaton ~S, in
which a set S' of final states has been indicated, while
the alphabet B is empty (i.e. XeS XA X S), represents
the event L~ consisting of all words a I ... an EA * such
that there are states S 1> 'Sn + I for which S I E S I ,
Sn+IES' and (Siaisi+I)EX is true for any i=l, ... ,no
The class of events that can be represented by the auto-
maton ~S, is identical with the class of regular events,
i.e. with respect to these aspects of their behaviour
non-deterministic automata are equivalent to finite
automata. However, the highly general nature of the
concept of a non-deterministic automaton is reflected
in the fact that a different number of states may be
necessary to represent the same event using a non-
deterministic automaton and using a finite automaton.
There are events which are representable by a non-
299
AUTOMATON
deterministic automaton with m states and by a finite
automaton with 2m states, but not by any finite auto-
maton with a smaller number of states.
A special subclass of non-deterministic automata is
formed by the so-called partial automata, in which the
relation X is a partial function mapping the set S XA
into S X B, and which realize partial finitely-determined
functions.
The term 'asynchronous automaton' usually denotes
one of the two following types of automata. The first
type comprises automata of type (A, S, B, cp, l/J), in
which the output function l/J maps the set S XA into B*
(for a finite automaton l/J maps S XA into B). These
automata are mainly used in coding theory. The second
type comprises finite automata whose transition func-
tion cp has the following property: cp(s, aa)=cp(s, a) for
all s and a. These automata are used in coding theory
and also in modelling certain systems in technology
and biology.
3. Automata with a variable structure (third group) are
finite automata 2f=(A XA, S, B, cp, l/J) with two input
channels, together with some fixed infinite sequence a
(a superword) in the alphabet A. Arbitrary words in
the alphabet A are sent to the first input (channel) of
such an automaton, while initial segments of the same
length of the sequence a are sent to the second input
(channel). This imposes a restriction on the set of pairs
of input words. If the automaton with a variable struc-
ture is considered as an automaton having only the first
input (into which any word of A can be fed), its transi-
tion and its output functions will depend on the length
of the input word which was fed in. In its behaviour,
an automaton with a variable structure is equivalent to
an infinite automaton with finite input and output
alphabets and with a countable set of states.
4. Fuzzy automata constitute a generalization of the
concept of a finite automaton, obtained by replacing
the transition and output functions by fuzzy relations.
A fuzzy subset of a set M is defined as a function which
maps M into the segment [0, 1]. Accordingly, in a fuzzy
automaton, the transition and the output functions are
replaced by functions mapping the sets S XA X Sand
S X A X B into [0, 1], where S is the set of states, A is
the input alphabet and B is the output alphabet. The
basic concepts and problems typical of finite automata
have natur:ll generalizatilJ!1s to fll7zy automata: this
applies, in particular, to the representation of fuzzy
events and the realization of fuzzy relations. Fuzzy
automata are mathematical models for certain recogni-
tion mechanisms, and are used in pattern recognition.
5. Special classes of finite automata.
Automata without memory are single-state finite auto-
mata or automata equivalent to them. In such auto-
mata each output letter is fully determined by the input
300
letter which is fed in at that moment. Such automata
are frequently referred to as functional elements and are
extensively used in problems of synthesis of automata.
Automata with a finite storage space (or automata with
a finite memory) are finite automata in which each out-
put letter is fully determined by a bounded segment of
the input word fed in during the preceding moments of
time, irrespective of the initial state. For automata with
a finite memory and with n states the length of such a
segment does not exceed n (n - I) / 2, and this max-
imum value is in fact attained for some automata.
Automata with a finite memory are called self-
regulating if, at a given time t, the output letter at any
moment 7"~t is independent of the initial state. Such
automata are used in coding theory (cf. Coding and
decoding), and the modification of the automata usually
employed in such cases satisfies the above condition
not for the set of all input words, but only for some
subset of this set. Automata with a finite memory are
realized by logical networks without feedback.
Reverse automata or automata without loss of informa-
tion are finite automata which realize one-to-one func-
tions. Such automata are also used in coding theory.
Micro-approach. There are three types of generalized
structural automata, which may be considered as gen-
eralized logical networks: 1) generalized logical net-
works with a permanent structure, in which both the
elements and the relation between them remain
unchanged when the automaton is functioning; 2) gen-
eralized logical networks with a variable structure; and
3) generalized logical networks made of volume ele-
ments and connections. The main classes of such auto-
mata are described below.
1. Generalized logical networks with a permanent struc-
ture. These include mosaic structures and iterative net-
works, which are finite fragments of mosaic structures
with a similar range of problems.
Mosaic structures are infinite unions of transition sys-
tems (A, S, cp) (i.e. of finite automata of the type
(A, S, S, cp, cp), where the output function is identical
with the transition function and the output alphabet is
identical with the set of states). Such systems are placed
at points with integer coordinates (integral points) of
the n-dimensional space, while for each such point
there is defined a finite set of integral points, called its
neighbourhood. The input alphabet of the transition
system placed at a certain point is the Cartesian pro-
duct of the sets of states of the transition systems
placed at the points in its neighbourhood.
A mosaic structure may be regarded as an infinite
automaton whose input and output alphabets as well as
its set of states are equal and are the infinite Cartesian
product of the sets of states of all the transition systems
contained in it. This makes it possible to reduce many

problems for mosaic structures to problems for infinite
automata. Problems specific to mosaic structures
include modelling of effective procedures; in particular,
of computational processes. Mosaic structures in which
arbitrary automata instead of transition systems are
used are also occasionally considered.
Uniform structures form an important class of mosaic
structures. If all the transition systems are identical
and if the neighbourhood of any point is obtained by a
parallel translation of a certain fixed neighbourhood,
the mosaic structure is known as a uniform structure or
a cellular automaton. It is usually assumed, in this con-
text, that there is some 'stable' state of the transitional
system, which is preserved if the input word is a tuple
each term of which corresponds to this state. Typical
problems on uniform structures are problems of self-
reproduction and the 'Garden-of-Eden' problem.
At any moment the states of the transition systems
located at points in an integral lattice generate a kind
of spatial mosaic pattern, which is usually denoted as a
configuration. A configuration containing only a finite
set of transition systems whose states are unstable (the
excited part) is called finite. The self-reproduction prob-
lem consists of finding out whether finite configurations
exist that, during the operation of the uniform system,
pass into configurations, whose excited part is the mul-
tiply iterated excited part of the original configuration.
The 'Garden oj Eden' is a term denoting a configuration
which cannot result from configurations different from
itself. The 'Garden-oj-Eden problem' is to determine the
existence of 'Gardens of Eden' for a given uniform
structure.
2. Generalized logical networks with a variable struc-
ture. There exist various types of such logical net-
works. The most general ones include mosaic structures
in which both the neighbourhoods of the elements and
the elements themselves vary during the operation. As
an example of such a generalized logical network one
may mention the one-dimensional structure simulating
the functioning of a Turing machine with input. In this
case one of the nodes of the one-dimensional network
corresponds to the control mechanism, while the other
nodes correspond to the tape cells, which are con-
sidered as transition systems in which the letters of the
working alphabet of the Turing machine serve as the
input letters and states. The commutation is determined
by the position of the head.
Another type of generalized logical networks with a
variable structure are the so-called growing automata.
They are uniform structures whose excited part
increases during the operation. There are a number of
models of such automata, which simulate various
features of real mechanisms.
3. Generalized logical networks oj volume elements are
AUTOMATON. BEHAVIOUR OF AN
distinguished by the fact that a certain volume is
assigned to their elementary automata and their mutual
connections. As a result there arises the problem of
synthesis of automata with minimum possible volume.
The term 'automaton' is also used in concepts such
as two-sided, multi-tape, multi-headed, linearly
bounded, etc., automata, which are in fact modifica-
tions of the Turing machine. Sometimes all abstract
machines are included in the concept of automata.
References
[I] McCullOCH, W.S. and PIrrs, W.: 'A logical calculus of ideas
immanent in nervous activity', Bull. Math. Biophys. S (1943),
115-133.
[2] KLEENE, S.c.: 'Reproduction of events in nerve sets and finite
automata', in Automata studies, Princeton Univ. Press, 1956,
pp. 3-14.
[3] BURKS, A.W. and WRIGHT, J.B.: 'Sequence generators, graphs,
and forma11anguages', Inform. and ControlS (1962), 204-212.
[4A] GLUSHIKOV, V.M.: 'The abstract theory of automata', Russian
Math. Surveys 16, no. 5 (1961), 1-53. (Uspekhi Mat. Nauk 16,
no. 5 (1961), 3-62)
[4B] GWSHIKOV, V.M.: Uspekhi Mat. Nauk 17, no. 2 (1962), 270.
[5] RABIN, M.O. and SCOTT, D.: 'Finite automata and their deci-
sion problems', IBM J. Res. Develop. 3 (1959), 114-125.
[6] ZADEH, L.A.: 'Probability measures of fuzzy events', J. Math.
Anal. Appl. 23 (1968),421-427.
[7] ALAD'EV, V.Z.: On the theory of uniform structures, Tallinn,
1972 (in Russian).
V.B. Kudryavtsev
Yu.I. Yanov
Editorial comments. See the Editorial Comment to the
article Automata, equivalence of.
AMS 1980 Subject Classification: 03005, 18820,
680XX, 93AXX
AUTOMATON, BEHAVIOUR OF AN - A mathemat-
ical concept describing the interaction of the automaton
with its external environment. Thus, for a finite auto-
maton (cf. Automaton, fmite) the external environment
usually is the set of input words, while its behaviour is
the word function generated by the automaton or the
event (sometimes the super-event) represented by it.
For an automaton over terms (see Automata, algebraic
theory of), the set of constant terms is the environment;
its behaviour is the class of terms whose values, calcu-
lated using the automaton, belong to the isolated subset
of elements of the corresponding algebra. For a mosaic
structure the environment is the set of initial configura-
tions, while its behaviour consists of the sequences of
configurations generated at the various moments of the
time set. In general, the concept of the behaviour of
(more general) automata is similar to, or is essentially a
slight modification of, the concept of the behaviour of
finite automata.
The so-called behaviour oj an automaton in a random
medium is a special case. Here, 'medium' is understood
to mean a probabilistic automaton ~, which transforms
the output signals of the automaton m: under considera-
301
AUTOMATON, BEHAVIOUR OF AN
tion to input signals for 2. Thus, it may be assumed
that the automaton 2 in the random medium ~ is an
autonomous logical network, constructed out of the
automata 2 and ~ by connecting the output of each
automaton with the input of the other. The behaviour
of the automaton 2 in the random medium ~ may
then be regarded as the functioning of this autonomous
logical network. The medium ~ is said to be stationary
if it is a single-state automaton. If the output signals of
the automaton ~ are regarded as 'rewards' or 'punish-
ments' for the automaton 2, the natural problem arises
of how to construct an automaton 2 whose behaviour
in the medium ~ would be optimal, i.e. would yield the
highest possible gain in that medium. It is usually
assumed that the output alphabet of the medium ~
consists of the letters 0 and 1 and that, as a response to
the output signals b \, ... ,bk of the automaton 2, the
letter 1 is put out, respectively, with probabilities
p J, . . ,Pk. Only the letter 1 is considered to be the
'reward' of the automaton 2.
If the medium ~ is stationary, the set of states of the
autonomous logical network becomes identical with the
set of states of the automaton 2. If, in addition, the
output letter of the automaton 2 is unambiguously
determined by its state, the functioning of this logical
network may be described by a stochastic matrix Q of
state transitions. The case which is usually considered is
that of an ergodic matrix (cf. Ergodicity). Then the
function
k k k
W(~,)8) = ~aiPi - ~ai(I-Pi) = ~ai(2pi -1)
i=1 i=1 i=1
is defined, where (Ji is the sum of the final probabilities
of all the states which determine the output letter hi.
Then
min(2p, -1) ,;;:;; W(~,)8) ,;;:;; max(2pi -1).
, ,
If the output signals of the automaton 2 do not depend
on the effects of the medium and if they are equally
probable, i.e. if (J\ = ... = (Jk = 1 / k, one has
W(~,)8) = Wo ?'-1.
2 k
= ki
The function W(2,~) is the mathematical expectation
of a variable called the gain of the automaton 2 in the
medium ~. One says that the behaviour of the automa-
ton 2r til a medium IR is goal-conforming if
W(2, ~ Woo The problem of optimal behaviour in a
random medium can now be stated as follows. One has
to construct a so-called asymptotically optimal sequence
of automata 21'\, 21'2, ... , such that the mathematical
expectation of the gain of automaton 2ln tends to the
maximum possible gain in this medium, which is
d max = maxi {2Pi -1: i = 1, ... ,k} as n increases. In the
case here considered such a sequence is formed by so-
called automata with linear tactics under the condition
302
that d rnax :;;;;.1 /2. (An automaton 2n with a linear tac-
tic, with an output alphabet of k letters, has kn states
sij' i = 1, ... ,k, j = I, ... ,n, and the following transi-
tion function cp and output function 0/:
q,(sij' 1) = Si,j+ I> if J = 1, ... ,n -1,
q,(Sin, 1) = Sin, q,(sij'O) = Si,j-I> if J=2, ... ,n,
q,(Sjj, 0) = SHI,1> !J,(sij,a) = hi,
j = 1, ... , k, J = 1, ... , n.)
The rules governing asymptotically optimal behaviour
in a stationary random medium were first studied in
mathematical statistics. However, the results obtained
in that field can obviously be translated into the
language of the theory of automata.
The behaviour of automata is also studied in more
complex media, as is the behaviour of collectives of
automata in random media. In the latter case the auto-
mata are considered to be players, and the rules of the
game played by the automata are considered as the
medium.
References
[I) BARZDIN', Y A.M. and 'fRAKlrrENBROT, B.A.: Finite automata:
behavior and syntheSiS, North-Holland, 1973 (translated from
the Russian).
[2) TSETLIN, M.L.: Studies on the theory of automata and modelling
of biological systems, Moscow, 1969 (in Russian).
[3) ROBBINS, H.: 'Sequential decision problems with a finite
memory', Proc. Nat. Acad. Sci. USA 42, no. 12 (1956), 920-923.
VB. Kudryavtsev
Yu.I. Yanov
Editorial comments. See the Editorial Comment to the
article Automata, equivalence of.
AMS 1980 Subject Classification: 03005, 680XX
AUTOMATON, FINITE - A mathematical model of
a system with a finite memory which processes discrete
information. Finite automata are one of the most
important types of control systems (cf. Control sys-
tem). Substantially, a finite automaton may be
described as a system with input and output channels
that is in one of n states s \, ... ,Sn at any discrete
moment of time. These moments constitute the time set.
At each such moment signals - that is, letters from the
input alphabet - are fed into the input channel and
signals - that is, letters from the output alphabet -
are produced by the output channel. Depending on the
particular poini uf view, :;uch system may include for-
mal systems (d. Formal system), real automata, living
organisms, etc.
The concept of a finite automaton may be defined
from different points of view. In the macro approach,
i.e. when the only feature of interest is the external
behaviour of the system, a finite automaton may be
defined as a class of functions, as a finite oriented
graph or (in algebraic form) as a finite algebra with

unary operations (cf. Automata, methods of specification
of). If the micro-approach is adopted, a finite automaton
is defined as a set of elements and a scheme of their
interconnections, i.e. not only the functioning of the
automaton but also its structure is considered. This
conception is accordingly known as structural, while the
finite automata themselves are known as structural
automata or automata networks.
Macro approach. A finite automaton is a system
(A, S, B, cp,l/I) where A, S, B are finite alphabets, usu-
ally non-empty, known, respectively, as the input alpha-
bet, the set of states and the output alphabet; cp is the
transition function, which maps the set S XA into S; 1/1
is the output function, which maps S XA into B. Such
finite automata are sometimes known as Mealy auto-
mata. If the output function 1/1 maps S into B (i.e. is
independent of the letters of the input alphabet), the
finite automaton is known as a Moore automaton. Any
Moore automaton is also a Mealy automaton.
The most important characteristic of a finite automa-
ton is its behaviour (cf. Automaton, behaviour of an),
which may be defined in different ways. Depending on
the kind of behaviour under consideration, finite auto-
mata may be grouped into transformers, acceptors
(identifiers), generators, etc. In order to define the main
types of behaviour of finite automata, the functions cp
and 1/1 are extended to the set S XA' (where A' is the
set of all words over A, including the empty word !\):
cp(s, !\) = s, cp(s, aa) = cp(cp(s, a), a),
Ij;(s, !\) = !\, Ij;(s, aa) = o/(cp(s, a), a),
where s ES, aEA', a EA, while aa denotes the word
obtained by adjoining the letter a to the word a. Thus,
the extension of the functions cp(s, a) and I/J(s, a) to
arbitrary s and a describes, respectively, the state into
which the automaton passes from the state s under the
effect of the input word a, and the output letter which
is produced by the automaton at the moment of entry
of the last letter of the input word a. Let al n denote the
initia! segment of length n of the word a, and let ~s, a)
and I/J(s, a) be the words over S and B defined as fol-
lows:
cp(s, a) = cp(s, all )lj;(s, ah) ... cp(s, a),
Ij;(s, a) = Ij;(s, all )lj;(s, ah) ... Ij;(s, a).
The functions ~s, a) and ~s, a) describe the sequence
of states assumed by the automaton during inputting
the letters of the word a, and also the output word, i.e.
the sequence of letters of the output alphabet put out
by the automaton when acted upon by the input word
a. The ternary relation
F = {(a,~s, a),J,(s, a)): aEA', SES}
is known as the functioning (operation) of the finite
automaton ~=(A, S, B, CP,I/I). The functions ~ and ~
AUTOMATON. FINITE
are extended in a natural manner to infinite sequences
(superwords) over A. For this reason, the functioning of
a finite automaton is sometimes understood to mean a
relation of the type F in which a is an arbitrary super-
word. In such a case the value of the function cp(s, a) is
defined as the set of those and only those states which
are encountered in the sequence ~s, a) an infinite
number of times.
A finite automaton with a designated initial state s 1
is known as initialized and is denoted by ~Sl' The
behaviour of an initialized finite automaton ~SI is usually
defined as one of the following four relations.
1) The function f(a)=~s], a), which maps A' into
B' (or A 00 into Boo, where A 00 and Boo are the sets of
all superwords over A and B respectively). This func-
tion is called computable or realizable by the initialized
finite automaton ~SI
2) The set Ls' CA', which is defined for a given
S' C S of the final states as follows:
Ls' = {aEA': cp(sl,a)ES'}.
The set Ls is called an event representable by the finite
automaton ~SI with set S' of final states.
3) The set of values of the function f(a)=~(s 1, a) for
all possible a from A " called the set enumerable by the
given finite automaton ~SI
4) The set M'Y CA 00, defined for a system y of sub-
sets of S as follows:
My = {aEA 00: CP(SI, a)EY}.
The set M'Y is known as a super-event, representable by
the finite automaton ~SI with a system y of subsets of
final states. Finite automata with behaviour of type 1)
are known as finite transformers, while those with
behaviour of type 2) are known as finite identifiers or
acceptors.
If the Cartesian products A 1 X ... XAm and
B 1 X ... X Bn are taken as the input and output
alphabets, respectively, behaviour of type 1) will be a
tuple of n functions of m arguments. One says in such
a case that the automaton has m inputs and n outputs,
the alphabets AI, ... ,Am and B 1, . . . ,Bn being called
input and output alphabets of such an automaton,
respectively. The class of events representable by a fin-
ite automaton and the class of functions computable by
a finite automaton may be described by various
mathematical tools. The principal result is that events
representable by a finite automaton are identical with
so-called regular events, while functions computable by
a finite automaton are identical with finitely-
determinded functions (d. Finitely-determinded func-
tion). In addition, the class of sets enumerable by a fin-
ite automaton is identical with the class of events
representable by a finite automaton.
303
AUTOMATON, FINITE
Regarding behaviour of types 2), 3) and 4), Moore
automata are equivalent to Mealy automata in the
sense that to each Mealy automaton there corresponds
an equivalent (i.e. displaying the same behaviour)
Moore automaton, and vice versa. This is not true of
behaviour of type 1). However, if in a Moore automa-
ton functions of the form
;j{s, a) = t[,{t/>(s, all )N{t/>(s, ah ... t[,{t/>(s, a,
are taken instead of ~, Moore automata will be
equivalent to Mealy automata.
A Moore automaton m=(A, S, B, cp,1/;) is said to be
an autonomous automaton, or a generator, if the transi-
tion function is independent of the letters of the input
alphabet. The behaviour of an initialized autonomous
automaton m SI is usually understood to mean the super-
word
where cpn+l(s)=cp(cpn(s)). This sequence IS periodic
apart from a certain initial segment.
A finite automaton mis said to be a transition system
if B = S and if, for any s from S, the equation
~s, a)=s is true, or if the output alphabet B is empty.
Thus, a transition system is completely defined by the
three parameters (A, S, cp).
The concepts of a finite automaton and functioning
of a finite automaton may be defined by several
equivalent methods (cf. Automata, methods of specifica-
tion of; Automaton, behaviour of an). The so-called
canonical equations are extensively employed. For any
word a let aCt) denote the t-th letter of a and let I a I
denote the length of a. The functioning F of a finite
automaton mwill then comprise those and only those
word triplets (a, a, 13) that satisfy the following condi-
tions: 1) I a I = I 13 I = I a I - 1; 2) for each t such that
1~t ~ I a I the equalities a(t + 1) = cp(a(t), aCt)),
f3(t)=t/;(a(t), aCt)) hold. To define the behaviour of an
initialized finite automaton mSI one must add the
equality a(l) =s I. The totality of these equalities unam-
biguously determines the behaviour of an initialized
finite automaton, and is usually referred to as the
canonical equations. Canonical equations are exten-
sively employed in the analysis and synthesis of auto-
mata.
Micro-approach.
l--r.-:.l-----t~:-r--l
____~ __ J
l-~T--
-t~-:-t -r.~~t
Fig. 1.
Consider a set of elements consisting of a finite auto-
304
maton as defined above with input and output alpha-
bets of the form An, where A is a finite alphabet, ident-
ical for all elements. The construction rules of struc-
tural finite automata from these elements describe the
permitted unions (identifications) of the inputs and out-
puts, and thus also define the input and output sets of
the finite automaton thus obtained.
I
g
... t
r--b--i--l
! t... i
I
... J
L_______
Fig. 2.
The most important example of such finite automata
are logical networks. One version of this concept is
presented below. In this case A = {a, I} and the ele-
ments are the so-called functional elements, which
represent a finite automaton with a single state, and
also certain special Moore automata, called delays.
These are distinguished by the fact that they have two
states, which are usually denoted by the letters and 1
of the input alphabet, and that the transition and out-
put functions satisfy the conditions cp(a, b)=b and
t/;( a) = a. Since functional elements have only one state,
they are completely determined by the output function
1/;, which in this case is a Boolean function of n argu-
ments, where n is the number of inputs of the func-
tional elements. The elements themselves are initialized
logical networks whose inputs and outputs correspond,
respectively, to the inputs and outputs of the elements.
Subsequent construction of more complex logical net-
works is performed according to the following rules.
1. The union of two logical networks is a logical net-
work whose inputs and outputs are, respectively, those
of the two logical networks (Fig. 1).
2. The result of identifying two arbitrary inputs of a
logical network is a logical network whose outputs are
the same as those of the initial logical network and
whose inputs are the same as those of the initial logical
network, except for one of the identified inputs (Fig. 2).
3. The result of connecting one output of a logical
network to an input of another logical network IS a log-
ical network. Its inputs are all the inputs of the first
logical network and those inputs of the second logical
network that are not identified with the output of the
first logical network; the outputs are all the outputs of
both logical networks (Fig. 3).
4. The result of identification of the output of a logi-
cal network that is the output of a delay element with
any input of the same logical network is a logical net-

Fig. 3.
work whose inputs are all the inputs of the initial logi-
cal network except for the identified input, and whose
outputs are all the outputs of the initial logical network
(Fig. 4). With certain restrictions this rule may also be
applied to outputs that are not outputs of a delay ele-
ment of a logical network.
i40
r--~-------1
I f ~ I
I 'I i
l--i-.-::- r ___ J
Fig. 4.
5. In any logical network it is possible to consider as
output only some of the outputs as defined in 1 - 4
above. Logical networks obtained from functional ele-
ments by the rules 1, 2, 3 and 5 are usually said to be
diagrams of functional elements.
The essential functioning of a logical network may be
substantially described as follows. Let specified input
letters be assigned to all inputs of the logical network
at the moment t and let the states of the delay elements
be given. In this situation letters will be assigned to all
inputs and outputs of the elements of the logical net-
work, in particular to all outputs of the logical network,
in accordance with the following rules. If the letters
ail' ... ,ai. have been assigned to the inputs of a func-
tional element with output function tf;(x1' ... ,xn ), the
letter assigned to its output at this moment will be the
value tf;(ail, ... ,aiJ. If at the moment t the delay is in
state a, the letter a is assigned at the same moment to
the output. Identical letters are assigned to identified
inputs and outputs. Moreover, at the moment t + 1 the
states of the delay are determined by the input letters
at the moment t, as was stated above, i.e. cf>(a, b)=b.
Thus, if the initial states of the delay are specified, the
logical network defines a certain mapping of input
sequences over the alphabet A m into output sequences
over the alphabet An, where m ~ 1 is the number of
inputs and n ~ 1 is the number of outputs of the logical
network. The class of such mappings is identical with
the class of functions realized by finite automata in the
first meaning of the word (i.e. with the class of
finitely-determined functions), since the functioning of
AUTOMATON. FINITE
the logical network as described above is identical with
the functioning of the finite automaton
(A m, S, An, cp, tV), where m is the number of inputs, n is
the number of outputs of the logical network and S is
the Cartesian product of the sets of states of all delays
in the logical network; the transition function cp is a
coordinate-wise application of the transition functions
of the delays, while the output function tV is determined
in accordance with the functioning of functional ele-
ments and delay elements as described above.
For instance, let the elements consist of delays
(marked as rectangles in Fig. 5) and of functional ele-
ments with output functions X, x Vy and x&y (marked
as triangles with the respective symbols of the func-
tions).
Fig. 5.
Fig. 5 shows a logical network which, at the moment of
time t, outputs letter 1 if and only if the input letters at
the moments 0, ... ,t contain an odd number of letters
'1' (at the initial moment the delay has the value 0;
letters a and b denote, respectively, the input and the
output of the logical network). If s(t), a(t) and b(t)
denote, respectively, the state, the input letter and the
output letter at the moment t, the functioning of such a
logical network may be defined by the canonical equa-
tions:
s(O) = 0, s(t + 1) = s(t)+a(t) (mod 2), b(t) = s(t).
If the macro-approach is adopted, this automaton can
be represented in the form (A, S, A, cp, tV) where
A =S={O, I}, cf>(s, a)=s+a (mod 2) and tf;(s, a)=s.
The concept of a finite automaton was the starting
point of the modern theory of automata (cf. Automata,
theory 01), which is also concerned with various modifi-
cations and generalizations of this concept.
References
[1] SHANNON, C. and MACCARTHY, J. (EDs.): Automata studies,
Princeton Univ. Press, 1956.
[2] GLUSHKOV, Y.M.: Synthesis of numerical automata, Moscow,
1962 (in Russian).
[3] KORBINSKII, N.E. and TRAlrnTENBROT, B.A.: Introduction to
305
AUTOMATON, FINITE
the theory offinite automata, Moscow, 1962 (in Russian).
V.B. Kudryavtsev
Yu.I. Yanov
Editorial comments. See the Editorial Comments to the
article Automata, equivalence of.
AMS 1980 Subject Classification: 03D05, 68DXX
AUTOMATON, PROBABILISTIC, stochastic automa-
ton - A generalization of a finite automaton (cf. Auto-
maton, rmite) in which the transition and the output
functions are random functions. In other words, a pro-
babilistic automaton may be defined as a system
(A, S, B, 1/>,1/1) where A, S, B are finite alphabets having
the same meaning as for finite automata, while 1/>, 1/1 are
random functions mapping S XA into S and B, respec-
tively, and are represented by systems of probability
measures I/>s,a, 1/Is,a, defined for all a from A and s from
S, on S and B, respectively. These measures are usually
given by a stochastic matrix (cf. Automata, methods of
specification of). If these probability measures can
assume only two values and 1, the concept of a pro-
babilistic automaton is practically identical with that of
a deterministic automaton. Autonomous probabilistic
automata without output are essentially equivalent to
discrete Markov chains. The functioning of a proba-
bilistic automaton is determined similarly to that of a
non-deterministic automaton, the initial state being
determined by giving a probability measure (J on S. If
the probabilistic automaton is in state s with a certain
probability p and accepts the input letter a, it will, with
probability p-w(s, a, s', b) pass into state s', and will
output the letter b of the output alphabet.
As with finite automata, probabilistic automata can
be classified into transformers and acceptors, in accor-
dance with their type of behaviour. In the former case
the automaton operates by transforming input words
into output words and into words in the alphabet of
states with certain probabilities. If the words are of the
same length, these probabilities form a probability
measure, since such behaviour may be regarded as giv-
ing a countable system of such measures. In the latter
case, a subset S' \: S of final states and a number II.
from [0, 1], called the cut point, are given. An event
which is representable by the probabilistic acceptor
2!" = (A, S, 1/>, S', 11.), where I/> is a random function
1I1aWllig S XA into S a.nd is defined by a system of
probability measures I/>s,a defined on S, consists of all
words over A under the effect of which the automaton
passes into one of the final states with probability at
least 11.. As distinct from finite automata, probabilistic
automata can represent a continuum class of events.
Moreover, even a single probabilistic automaton may
represent a continuum class of events if ,\ is varied. If
the input alphabet is a one-letter alphabet, each proba-
306
bilistic automaton merely represents a countable class
of events, which also includes irregular events in the
general case. At special cut points, the so-called iso-
lated points, the probabilistic automaton represents
only regular events. A number II. from the segment
[0, 1] is known as an isolated cut point for this automa-
ton if there exists a positive number 8 such that for any
input word the probability of transforming the proba-
bilistic automaton by this word into its final state
differs by at least 8 from 11..
Most of the concepts and problems typical of finite
automata may be generalized, in different variants, to
probabilistic automata. Many of the concepts and
problems preserve the properties of finite automata.
For instance, it is possible to introduce the concept of
equivalence of states so as to retain the well-known
theorem on the possibility of differentiating between
states by means of a simple experiment (cf. Automata,
experiments with). On the other hand, unlike finite
automata, whose minimal form is unambiguously
defined (up to an isomorphism), a continuum of
equivalent minimal probabilistic automata may exist
for a given probabilistic automaton.
There are different ways and means of setting up
probabilistic automata. For instance, a probabilistic
automaton may be represented by a deterministic auto-
maton with two inputs into one of which is fed a ran-
dom sequence of input letters. Probabilistic automata
are mathematical models of many real mechanisms and
are used in studying the behaviour of living organisms.
References
[1] BUKHAREV, R.O.: Probabilistic automata, Kazan, 1970 (in Rus-
sian).
[2] STARKE, P.: Abstrakte Automaten, Deutsch. Verlag Wissen-
schaft., 1969.
V.B. Kudryavtsev
Yu.I. Yanov
AMS 1980 Subject Classification: 68025
AUTOMORPHIC FORM - A meromorphic function
on a bounded domain D of the complex space en that,
for some discrete group r operating on this domain,
satisfies an equation:
f(y(x = F;m(x)f(x), XED, yEf,
Here j y(x) is the Jacobian of the mapping x ~y(x) and
m is an integer known as the weight of the automorphic
form. If the group r acts fixed-point free, then auto-
morphic forms define differential forms on the quotient
space D / r and vice versa. Automorphic forms may be
used in the construction of non-trivial automorphic
functions (cf. Automorphic function). It has been
shown that if g (x) is a function that is holomorphic
and bounded on a domain D, then the series
2:g(y(xj~' (x)
yE"i'

converges for large values of m, thus representing a
non-trivial automorphic function of weight m. These
series are called Poincare theta-series.
The classical definition of automorphic forms, given
above, has recently served as the starting point of a
far-reaching generalization in the theory of discrete
subgroups of Lie groups and adele groups.
References
[1] POINCARE, H.: Oeuvres de Henri Poincare, Gauthier-Villars,
1916-1965.
[2] SIEGEL, c.L.: Automorphe Funktionen in mehrerer Variablen,
Math. Inst. Gottingen, 1955.
A.N. Parshin
Editorial comments. References [A2] and [A3] can
serve to get some idea of modern developments and topics
in the theory of automorphic forms and its relations with
other parts of mathematics. (Cf. the comments to the article
Automorphic function for a more general notion).
References
[A1] BAILY, W.L. JR.: Introductory lectures on automorphic forms,
Iwanami Shoten & Princeton Univ. Press, 1973.
[A2] BOREL, A. and CASSELMAN, W. (EDS.): Automatic forms,
representations and L-functions, 1-2, Amer. Math. Soc.,
1979.
[A3] GELBART, S.S.: Automorphic forms on adele groups, Prince-
ton Univ. Press, 1975.
AMS 1980 Subject Classification: 10005, 32NXX
AUTOMORPHIC FUNCfION - A meromorphic func-
tion of several complex variables that is invariant under
some discrete group f of analytic transformations of a
given complex manifold M:
I(y(x)) = I(x), XEM, YEf.
Automorphic functions are often defined so as to
include only functions defined on a bounded connected
domain D of the n-dimensional complex space cn that
are invariant under a discrete group f of automor-
phisms of this domain.
The quotient space X =M / f can be given a com-
plex structure and automorphic functions are then
meromorphic functions on X. A large number of cases
studied concern a space X with a compactification X. It
is then natural to include in the definition of an auto-
morphic function the requirement that it can be
extended to the entire space X as a meromorphic func-
tion. If M =D (i.e. M is a bounded connected domain),
this condition must be required for n = I only (if n > I
or if M / f is compact, the condition is automatically
fulfilled). It can readily be shown that the automorphic
functions constitute a field K(f) and the study of this
field is one of the main tasks in the theory of auto-
morphic functions.
Automorphic functions of a single variable have been
very thoroughly studied. The theoretical foundations
were laid by F. Klein [1] and H. Poincare [2] in the
nineteenth century. The manifold M usually considered
AUTOMORPHIC FUNCTION
at that time is a simply-connected domain. Three cases
are distinguished: M=pl(C) (the complex projective
line, or the Riemann sphere), M=C and M=H (the
upper half-plane {ZEC: Imz>O}). In the first case the
discrete groups f are finite, the curves M / f are alge-
braic curves of genus 0 (cf. Genus of a curve) and, con-
sequently, the automorphic functions generate a field of
rational functions. Examples of automorphic functions
in the case M = C are periodic functions (thus, the
function e 2wiz is invariant under the translation group
{z~z +n: n EZ}) and, in particular, elliptic functions.
In this latter case, the curve M / f is compact and is
an elliptic curve, while the field K (f) is the field of all
algebraic functions on M / f. Finally, for M=H and
a discrete group f such that M / f is compact or has a
finite volume (in the Poincare metric), M / f is an alge-
braic curve and K(f) is again the field of algebraic
functions on M / r. The genus g of this curve may be
determined by constructing a fundamental domain for
f in the form of a polygon in the upper half-plane H
(here regarded as the Lobachevskii plane). The basic
method for constructing an automorphic function in
this situation is to consider the quotient of two auto-
morphic forms (cf. Automorphic form) of the same,
sufficiently large, weight. The method is due to Poin-
care, who used it to prove the results mentioned above
concerning the structure of the fields of automorphic
functions [2], [3], [4]. An analogous construction for
elliptic functions is to represent such functions in terms
of quotients of theta-functions. It can be shown, using
uniformization theory, that all fields of algebraic func-
tions of a single variable are obtained in this way [3J.
These results, which were obtained as early as the
nineteenth century, give a full description of the fields
of automorphic functions for n = I and of the groups f
such that the space H / f has finite volume. The case
of groups f for which H / f has infinite volume
(Kleinian groups) is much more difficult; the problems
involved are still being intensively investigated [5], [6].
In the twentieth century the theory of automorphic
functions concentrates on functions of several variables.
Perhaps the only case of automorphic functions of n
variables studied in detail in the nineteenth century
concerned Abelian functions, which are related to
Abelian varieties in a way similar to the relationship
between elliptic functions and elliptic curves [I], [7].
The first example of automorphic functions of n vari-
ables on a bounded domain D are the modular func-
tions of c.L. Siegel [7] (cf. Modular group). Their
domain of definition is an n-dimensional generalization
of the upper half-plane H, and is one of the main
examples of a bounded, symmetric domain. Siegel
must also be credited with the first general results
obtained concerning arbitrary automorphic functions
307
AUTOMORPHIC FUNCTION
on a bounded domain D. He generalized Poincare's
construction of automorphic functions, mentioned
above, and showed that the field K(f) always contains
at least n algebraically independent functions.
Subsequent efforts were aimed at exhibiting domains
D and groups f for which the following theorem on
algebraic relations is true. If JI, ... ,f" are algebraically
independent automorphic functions, then the field K(f)
is a finite algebraic extension of the field of rational
functions e(fJ, ... ,f,,) CK(f).
At the time of writing (1977) this theorem was
proved for the following cases: 1) if the quotient space
D / f is compact [7]; 2) if the group f is pseudo-
concave [8]; and 3) if D is a symmetric domain and f
is an arithmetic group. A pseudo-concave group is
defined as follows. Let X be a subdomain of a domain
D with closure also contained in D. In this situation a
boundary point Xo EaX is said to be pseudo-concave if
for any open neighbourhood U of Xo and for any func-
tion cp(x) regular in U there exists a point x E U n X
for which 1 cp(x) 1 ~ 1 cp(xo) I. A group f is said to be
pseudo-concave if there exists a subdomain XC D such
that each boundary point x EaX can be transformed by
means of an element of f into an interior point of X or
into a pseudo-concave point of the boundary ax.
The nature and the properties of the algebraic
varieties occurring in the theory of automorphic func-
tions of n variables have not been intensively studied,
as distinct from the case of a single variable.
Important generalizations of the concept of auto-
morphic functions - automorphic forms, theta-
functions (cf. Theta-function) and certain other gen-
eralizations - are all special cases of the following
general construction. Consider a fibre bundle (cf. Fibra-
tion) Lover M and an action of a group f on it. It is
then possible to consider the sections of L that are
invariant under r. An automorphic function is
obtained if the fibre bundle L and the action of the
group r are both trivial.
The study of automorphic functions revealed the
important role played by the group of automorphisms
of a domain D. It is in this way that the concepts and
methods of the theory of automorphic functions were
applied in the theory of algebraic groups, in which they
play an important part in the description of infinite-
di..~ensiona! representations [10]
From the very beginning of its development, the
theory of automorphic functions has been connected in
numerous ways with other branches of mathematics.
This applies in particular to algebraic geometry. In
addition to the results discussed above, methods in the
theory of automorphic functions are important in the
study of moduli varieties for objects such as algebraic
curves and Abelian varieties. Automorphic functions
30~
are also of importance in number theory. At the time
of writing they are the only tool in the study of zeta-
functions of algebraic varieties [ 11 ]. Another very
promising number-theoretical direction in the theory of
automorphic functions is the study of p-adic auto-
morphic functions and forms [9]. Finally, one must
mention the application of automorphic functions to
the study of ordinary differential equations in a com-
plex domain [I2] and in the construction of solutions of
algebraic equations of degrees higher than four.
References
[1] KLEIN, F.: Development of mathematics in the 19th century, 1,
Math. Sci. Press, 1979, Chapt.8 (translated from the Ger-
man).
[2] POINCARE, H.: Oeuvres de Henri Poincare, Vol. 4, Gauthier-
Villars, 1916-1965.
[3] FORD, L.R.: Automorphic functiOns, Chelsea, reprint, 1951.
[4] BATEMAN, H. and ERDELYI, A.: Higher transcendental functions.
Automorphicfunctions, Vol. 3, McGraw-Hill, 1955.
[5] La geometrie non-euclidienne dans la tMorie des fonctions auto-
morphes, Moscow, 1952 (in Russian; translated from the
French).
[6) KRA, I.: Automorphic forms and Kleinian groups, Benjamin,
1972.
[7) SIEGEL, C.L.: Automorphe Funktionen in mehrerer Variablen,
Math. lost. Gottingen, 1955.
[8) ANDREOTII, A. and GRAUERT, H.: 'Algebraische Korper von
automorphen Funktionen', Nachr. Akad. Wiss. G6ttingen 3
(1961).
[9) SERRE, J.-P., DELIGNE, P. and KUYK, W. (EDS.): Modular func-
tions of one variable. Proc. Internat. Summer Scool RUCA 1972,
1-3, Springer, 1973.
[10) JACQUET, H. and LANGLANDS, R.P.: Automorphicforms on
GL(2), Springer, 1970-1972.
[il) SHIMURA, G.: Introduction to the arithmetic theory of auto-
morphiC functiOns, Math. Soc. Japan, 1971.
[12) GoLUBEV, V.V.: Vorlesungen Uber Differentialgleichungen im
Komplexen, Deutsch. Verlag Wissenschaft., 1958 (translated
from the Russian).
A.N. Andrianov
A.N. Parshin
Editorial comments. The result mentioned above that
the field K(f) is a finite algebraic extension of a field of
rational functions C(f1, ... , fn) (the theorem on algebraic
relations) in the case of a symmetric domain 0 and an arith-
metic group r is due to, independently, W.L. Baily jr. and A.
Borel [A6], and 1.1. PyatetskiT-Shapiro [An
Let X be some kind of space (e.g. complex- or real-
analytic, a smooth manifold), r a group of automorphisms
of X and H a group acting on a space V. Let Mor(X, H) be
the set of morphisms from X into H. An automorphy factor
of r is a 1-cocycle (crossed homomorphism) j of r with
values in Mor(X, H). That means it is a mapping j: XX r--,>H
such that llX, yy') = j(x, y)j(xy, y"). All eXctli1pie is the Jaco-
bian of j as a diffeomorphism X--,>X (chain rule). An auto-
morphic form of type j is now a morphism f: X-'> V such that
f(x) = j(x, y) f(xy). Taking the Jacobian as an automorphy
factor and H=GL(C) acting on C via the m-th power of the
determinant one recovers the more classical notion of an
automorphic form of weight m, cf. Automorphic form. The
automorphy factor j can be used to define an action of r on
XX V by (x, v)y=(xy, j(x, y)v). If r now operates freely on

xx Vas a properly discontinuous group of transformations,
then (XX V) / f is a fibre bundle over X / f with fibre Vand
the automorphic forms are the cross-sections of this bundle,
or, equivalently, the f-equivariant cross-sections of the
trivial bundle XX V~x.
In a still more group-theoretical setting let G be a real
semi-simple Lie group with Lie algebra g. Identify the univer-
sal enveloping algebra Ug of g with the right-invariant dif-
ferential operators D( G) on G by extending the mapping
which assigns to a E g the corresponding right-invariant vec-
tor field. Let K be a maximal compact subgroup of G and f
a discrete subgroup and let p: K~GL(V) be a representation
of K. A smooth vector-valued function f: G~ V is called an
automorphic form for f if f(kgy) =p(k) f(g), (Zg)f is a finite
vector space, where Z(g)CUg=D(G) is the centre of Ug,
and f satisfies a certain growth condition. The link with the
notion 'automorphic form of type j' discussed just above is
provided by X = K \ G, the left coset space of Kin G, and a
canonical automorphy factor (with H= Kc) which can be
defined in this setting. Cf. [A 1] for more details on all this.
Besides the applications of automorphic functions in ordi-
nary differential equations and algebraic equations men-
tioned above there is also a most revealing connection
between the harmonic analysis of functions automorphic
with respect to a discrete subgroup of SL2 (R) and
Lax- Philips scattering theory applied to the non-Euclidean
wave equation, cf. [A4], [A5].
For more material closely related to automorphic forms
and automorphic functions, cf. also the articles Modular
form; Modular function; Fuchsian group; Discrete sub-
group; Discrete group of transformations.
References
[A1] BOREL, A.: 'Introduction to automorphic forms', in A. Borel
and G.D. Mostow (eds.): Algebraic groups and discontinuous
subgroups, Amer. Math. Soc., 1966, pp. 199-210.
[A2] FRICKE, R. and KLEIN, F.: Vorlesungen Dber die Theorie der
automorphen Funktionen, 1-2, Teubner, 1926.
[A3] BOREL, A. and CASSELMAN, W.: Automorphic forms,
representations and L-functions, 1-2, Amer. Math. Soc.,
1979.
[A4] FADDEEV, L.D. and PAVLOV, B.S.: 'Scattering theory and
automorphic functions', Proc. Steklov Inst. Math. 27 (1972),
161-198.
[AS] LAX, P.D. and PHILLIPS, R.S.: 'Scattering theory for auto-
morphic functions', BUll. Amer. Math. Soc. (New Ser.) 2
(1980), 261-296.
[A6] BAILY, W.L., JR. and BOREL, A.: 'Compactifications of arith-
metic quotients of bounded symmetric domains', Ann. of
Math. 84 (1966), 442-528.
[A7) PYATETSKlT-SHAPIRO, I.I.: 'Arithmetic groups on complex
domains', Russ. Math. Surveys 19 (1964),83-109.
AMS 1980 Subject Classification: 32NXX
AUTOMORPIDSM - An isomorphism (isomorphic
mapping) of a system of objects onto itself. The totality
of all automorphisms of an arbitrary algebraic system
forms a group, and the study of this group is an impor-
tant and powerful tool in the study of the properties of
the system itself (cf. Algebraic system, automorphism of an).
AMS 1980 Subject Classification: 18A20
ACT00iOMOUS SYSTEM
AUTONOMOUS SYSTEM of ordinary differential
equations - A system of ordinary differential equations
which does not explicitly contain the independent vari-
able t (time). The general form of a first-order auto-
nomous system in normal form is:
Xj = jj(Xh ... ,xn ), )=1, ... ,n,
or, in vector notation,
x = f(x). (1)
x
A non-autonomous system = f (t, x) can be reduced
to an autonomous one by introducing a new unknown
function Xn + I = t. Historically, autonomous systems
first appeared in descriptions of physical processes with
a finite number of degrees of freedom. They are also
called dynamical or conservative systems (cf. Dynamical
system).
A complex autonomous system of the form (1) is
equivalent to a real autonomous system with 2n
unknown functions
d d
dt (Rex) = Ref(x), dt (1m x) = Imf(x).
The essential contents of the theory of complex auto-
nomous systems - unlike in the real case - is found
in the case of an analytic f (x) (cf. Analytic theory of
differential equations).
Consider an analytic system with real coefficients
and its real solutions. Let x =q,(t) be an (arbitrary)
solution of the analytic system (1), let ~ =(L , t +) be
the interval in which it is defined, and let x (t; to, xo)
be the solution with initial data x 1/=10 =xo. Let G be a
domain in Rn and fEe 1(G). The point x EGis said to
be an equilibrium point, or a point of rest, of the auto-
nomous system (1) if f(xo)-O. The solution cp(t)=-X,
t E R = ( - 00, + 00), corresponds to such an equilibrium
point.
Local properties of solutions. I) If 4>(t) is a solution,
then 4>(t + c) is a solution for any c ER.
2) Existence: For any to ER, xO EG, a solution
x(t; to, xo) exists in a certain interval ~3t.
3) Smoothness: If fEcP(G), p~l, then
q,(t)E cP + 1 (~).
4) Dependence on parameters: Let f= f (x, a),
aEG a CRm , where G a is a domain; if fEcP(GXG a ),
p ;;:'1, then x (t; to, xO, a)E cP(~ X G a) (for more details
see [1] - [4]).
5) Let xO be a non-equilibrium point; then there
exist neighbourhoods V, W of the points xO,f(xo),
respectively, and a diffeomorphismy=h(x): V~W such
that the autonomous system has the form y =const in
W.
A substitution of variables x = <I>(y) in the auto-
nomous system (1) yields the system
(2)
309
AUTONOMOUS SYSTEM
where q; (y) is the Jacobi matrix.
Global properties of solutions. 1) Any solution x =4>(t)
of the autonomous system (1) may be extended to an
interval a=(I-, t +). If a=R, the solution is said to
be unboundedly extendable; if t + = + 00, t _ > - 00, the
solution is said to be unboundedly extendable forwards
in time (and, in a similar manner, backwards in time). If
t + < + 00 then, for any compact set K cO, xO EK,
there exists a T=7{K)<t + such that the point
x(t;to,xo) is outside K for t>7{K) (and, analogously,
for t _ > - 00; cf. Prolongation of solutions of differen-
tial equations).
2) The extension is unique in the sense that any two
solutions with common initial data are identical
throughout their range of definition.
3) Any solution of an autonomous system belongs to
one of the following three types: a) aperiodic, with
4>(t)=F4>(t2) for all t(=I=t2, tjER; b) periodic, non-
constant; or c) 4>(t)=const.
Geometric interpretation of an autonomous system. To
each solution x =4>(t) is assigned a coresponding curve
r: x =4>(t), tEa, inside the domain G. G is then said
to be the phase space of the autonomous system, r is
its trajectory in the phase space, and the solution is
interpreted as motion along this trajectory in the phase
space. The mapping gl: G~G defined by the formula
gIXO=x(t;O,xo) (Le. each point moves along the phase
trajectory during time t) is called the phase flow. In its
domain of definition the phase flow satisfies the follow-
ing conditions: 1) gl X is continuous in (t, x); and 2) the
II
group property +1, X =gIl gl, x.
The Liouville theorem is valid: Let D C G be a
domain with a finite volume and let VI be the volume of
the domain gl D C G, then
~: Lo = div f(x) dx.
(3)
For a Hamiltonian system, a consequence of (3) is
the conservation of the phase volume by the phase
flow. A second variant of (3) is obtained as follows. Let
x =cf>(t, a) be a family of solutions of (1),
a=(a), ... ,an-dEG a , let Ga be a domain and let
q,EC)(aXG a ), then
:r lnI(t, a) = div f(x), (3')
where I (t, a) =det ax / a(t, a).
Structure of phase trajectories. 1) Any two phase tra-
jectories have either no points in common or coincide.
2) Any phase trajectory belongs to one of the
following types: a) a smooth, simple, non-closed Jor-
dan arc; b) a cycle, i.e. a curve diffeomorphic to a cir-
cle; or c) a point (an equilibrium point). The local
structure of phase trajectories in a small neighbourhood
310
of a point other than an equilibrium point is trivial (cf.
local property 5) of the solutions): The family of phase
trajectories is diffeomorphic to a family of parallel
straight lines. For a linear autonomous system the
structure of phase trajectories in a neighbourhood of an
equilibrium point is known, since the autonomous sys-
tem is integrable [5J. For non-linear autonomous
systems this problem has not yet been completely
solved, even for n =2 (d. Qualitative theory of differen-
tial equations). One aspect of this problem is the ques-
tion of stability of an equilibrium point (cf. Stabllity
theory). A few results will be given below. Let xO,yO
be equilibrium points of the system (1), let
j = g(y) (1')
and let U, V be neighbourhoods of the points xO ,yo.
The systems (1) and (1') are said to be equivalent in
neighbourhoods of their equilibrium points xO ,yO if there
exist neighbourhoods U, V and a bijective mapping
h: U~V such that (hof)x=(g'oh)x (for XEU, fXEU,
(gl oh)x E V), i.e. as a result of the substitution y =h (x)
the trajectories of the autonomous system (1) go into
trajectories of the autonomous system (1'). The
equivalence is said to be differentiable (topolOgical) if h
is a diffeomorphism (homeomorphism). Let xO be an
equilibrium point of the autonomous system (1), let the
matrix / (xo) be non-degenerate, and let it not possess
any pure imaginary eigen values. Then the autonomous
system (1) in a neighbourhood of xO is topologically
equivalent to its linear part y =/ (xo)y. An important
example is the autonomous system x=Ax,y=By where
A, B are constant matrices with pure imaginary eigen
values and n > 2; it is not known when these auto-
nomous systems are topologically equivalent. One of
the most fundamental problems in the theory of auto-
nomous systems is that of the structure of the entire
family of phase trajectories. The most complete results
have been obtained for n = 2, but even in this case the
solution is far from complete.
References
[1) PETROWSKI, I.G. [I.G. PETROVSKIi): Vorlesungen aber die
Theorie der gew6hnlichen Differentialgleichungen, Teubner, 1954
(translated from the Russian).
(2) PONTRYAGIN, L.S.: Ordinary differential equations, Addison-
Wesley, 1962 (translated from the Russian).
(3) CODDINGTON, E.A. and LEVINSON, N.: Theory of ordinary dif-
ferential equations, McGraw-Hill, 1955.
(4) ARNOL'D, V.I.: Ordinary differential equations, M.LT., 1973
(translated from the Russian).
(5) NEMYTSKli, V.V. and STEPANOV, V.V.: Quantitative theory of
differential equations, Princeton Univ. Press, 1960 (translated
from the Russian).
M. V. Fedoryuk
Editorial comments. The question of the topological
equivalence (homeomorphic equivalence of orbit systems)
x y
of two linear dynamical systems = Ax and = By has been
examined and solved [A2]. Let V,t, VA, \.1 be the decompo-

sition of Rn corresponding to (generalized) eigen values of
A with positive, negative and zero real parts respectively,
and let AD be the restriction of A to \.1. Then the orbit sys-
tems of x=Ax and y=By are homeomorphic if
dim lit =dim vt, dim VA =dim Va and AD and BO are
linearly equivalent (i.e. AD = S-1 B OS for some constant
invertible matrix 5).
The analogous question for discrete dynamical systems:
when does there exist for two linear endomorph isms
A, B: Rn~Rn a homeomorphism t/> such that At/> = t/>B, is
much deeper (topological similarity of matrices) and has to
do with such things as lens spaces and the P.A. Smith con-
jecture. See [A1] for a full account. A theorem of G. de
Rham says that if the orthogonal matrices A and B are topo-
logically similar and the topological (non-linear) similarity
preserves the unit sphere and restricts to a diffeomorphism
in it, then the matrices A and B are also linearly similar.
Results of S.E. Cappell and J.L. Shaneson imply that for
matrices of eigen values of absolute value 1 (the crucial
case) topological and linear equivalence are the same for
dimensions ";;;5. Combined with results of Kuiper-Robbin
[A3] this gives a complete topological classification of
matrices of dimension at most 5. In higher dimensions it is
definitely not true in general that for matrices of eigen
values of modulus 1 topological similarity implies linear simi-
larity [A 1]. These results can be used to tell at what times t
the phase flows for continuous-time systems become topo-
logically equivalent. This can happen repeatedly even for
systems that are not topologically equivalent.
References
[A 1] CAPPELL, S.E. and SHANESON, J.L.: 'Non-linear similarity',
Ann. of Math. 113 (1981), 315-355.
[A2] KUIPER, N.H.: 'The topology of the solutions of a linear dif-
ferential equation on Rn', in Proc. Internat. Congress on
Manifolds Tokyo, 1973, pp. 195-203.
[A3] KUIPER, N.H. and ROBBIN, J.W.: 'Topological classification of
linear endomorphisms', Inv. Math. 19 (1973), 83-106.
AMS 1980 Subject Classification: 34A34
AUTONYMY - An expression (a word) used as its
own proper name. Such a use of an expression is said
to be autonymous (as distinct from its use in the ordi-
nary sense). For instance, if one says 'x forms part of
the equation x + 3 = 2', then one is using x as the name
of the letter x, and is using 'x + 3 = 2' as the name of
the expresion x+3=2. If one says '12 is divisible by
2', then the term '12' is used in a non-autonymous
manner (as the designation of an integer), while if one
says '12 consists of two digits' one says the term '12' is
used autonymously (as the name of itself).
In natural languages the context and the syntax are
usually sufficiently reliable indicators for distinguishing
between the autonymous and non-autonymous use of
an expression. However, such a differentiation may be
difficult in certain cases. Special care is then required
to avoid ambiguity: the object itself must be dis-
tinguished from its name (its denotation); a difference
AVERAGE
must be made between a given linguistic term used as a
name, and the term used for the object denoted by the
name. The differentiation between the denotation and
the object denoted may be achieved by using terms spe-
cially created for the purpose, or by using quotation
marks; the term between quotation marks is then
assumed to differ from the same term used without the
quotation marks. Autonymous use of a term merges the
meaning of the term with the term itself: it is both the
object itself and the name by which it is denoted.
References
[1] CHURCH, A.: Introduction to mathematical logic, 1, Princeton
Univ. Press, 1956.
[2] KLEENE, S.c.: Introduction to metamathematics, North-Holland,
1951.
[3] CURRY, H.B.: Foundations of mathematical logic, McGraw-Hill,
1963.
A.S. Kuzichev
AMS 1980 Subject Classification: 00A25
AVERAGE with weight q=(ql,'" ,qn), qi>O,
~qi = 1, of a set of real numbers a =(aj, ... ,an) - A
variable
ffiC</o(a, q) = t/>-I [~qjt/>(a,)],
I
where <t>(x) is a continous strictly-monotone function on
R. When <j>(x) = x', one obtains
ffiCr(a, q) = [~qjaT / r
and, in particular, when r = 1, qi = 1/ n, i = 1, ... , n,
1JR,(a, 1/ n)=m:(a) will be the arithmetic average of the
numbers aI, ... , an, while when r = - 1, it will be the
harmonic average. The concepts of the geometric average
(a) = (IIai)1 / n and the weighted geometric average
I
@(a, p) = [1/afT /~'p,
are introduced separately.
One of the basic results of the theory of averages is
the inequality (a)<m:(a), except when all ai are equal
to each other. Other results are:
1) IJRcp(ka,p)=klJRcp(a,p), k>O;
2) f1J4(a,p)=lJRcp(a,p) if and only if 1/;= acp +/3,
/3ER, a=O;
3) 1JR",(a,p)~lJRcp(a,p) if and only if cpol/;-I is a con-
vex function; in particular 1JR,(a,p)~IJRS<a,p) if r<s.
The concept of an average can be extended to infin-
ite sequences under the assumption that the
corresponding series and products converge, and to
other functions. The following is such an example:
b
t/>-I jp(x)<p(j(xdx
ffiC</o (t, p) = - - - " -a- b - - - - - - " ' -
jp(x)dx
a
311
AVERAGE
given the condition that J(x)~O almost everywhere on
the corresponding interval and that p(xO. Thus,
b b
jf(x)p(x)dx ",;;; lJJl.pif,p)jp(x)dx.
a a
References
[1) HARDY, G.H., LITTLEWOOD, J.E. and P6LYA, G.: Inequalities,
Cambridge Univ. Press, 1934.
VI. Sobolev
Editorial comments. Instead of 'average' the term
'mean' is also quite often used: arithmetic mean, geometric
mean, etc.
References
[A 1] MITRINOVlt, D.S.: Analytic inequalities, Springer, 1970.
[A2] MITRINOVlt, D.S.: Elementary inequalities, Noordhoff, 1964.
AMS 1980 Subject Classification: 26DXX
AVERAGE ROTATION of the argument argJ(t) of a
complex-valued, uniformly almost-periodic function
J (t)- A phenomenon consisting of the existence (given
certain conditions, see below) of the limit
lim arg f(t) = c.
' ..... 00 t (5) Doss, R.: 'On mean motion', Amer. J. Math. 79, no. 2 (1957),
The limit itself is also called the average rotation (mean
motion). If IJ (t) I~O for all t, then the selection of a
continuous branch of argJ (t) is implied. For analytic
almost-periodic functions, the concept of average rota-
tion can be retained even when J contains zeros.
Namely, the concepts of 'right' and 'left' arguments are
introduced, whose difference jumps by +k'TT at a zero
of multiplicity k of f, and one therefore also speaks of
right and left average rotation, unless they coincide, in
which case one simply speaks of average rotation.
The question of the average rotation arose in connec-
tion with the fact that in celestial mechanics the longi-
tude of the perihelion of a planet is expressed approxi-
mately as the argument of a certain trigonometric poly-
nomial
J.L. Lagrange studied two simple cases, namely, where
one of the I aj I is greater than the sum of the remain-
ing coefficients, and for n = 2; he noted that in other
cases the question is complicated. The study of this
question has been taken up only in the twentieth cen-
tury (for its history, see [I] - [3]). The final result that
a trigonOlllt:lllC polynomial ahvays has an average rot::l-
tion was stated in 1938 by B. Jessen (for its proof, see
[I)). (From the point of view of the theory of dynami-
cal systems, it is a matter of averaging a certain func-
tion on a torus along the trajectories of the flow
defined by shifts by elements of a one-parameter sub-
group. However, this function has singularities that
obstruct the automatic use of the corresponding general
theorem.) Even earlier than this, H. Bohr proved the
312
existence of an average rotation for any uniforml
almost-periodic function for which inf IJ (t) I >0 (see
[4]). In this case, the difference argJ(t)-ct is a uni-
formly almost-periodic function and is bounded. The
average rotation of analytic almost-periodic functions
in the general case has also been studied (see [1], [4]).
In this case, the average rotation does not always exist,
but if it does, then the difference argJ (t) - ct is not
necessarily bounded. Nonetheless, it can still possess
certain generalized properties of almost-periodicity; this
is true, in particular, for trigonometric polynomials [5].
Except in the analytic case, only isolated results con-
cerning the average rotation of a function J for which
IJ(t) I~O, inf IJ(t) I =0, exist (see [6], [7]).
References
[1) JESSEN, B. and TORNEHAVE, H.: 'Mean motions and zeros of
almost periodic functions', Acta Math. 77 (1945), 137-279.
[2) JESSEN, B.: 'Some aspects of the theory of almost periodic
functions', in Proc. Internat. Congress Mathematicians Amster-
dam, 1954, Vol. 1, North-Holland, 1954, pp. 304-351.
[3) WEYL, H.: 'Mean motion', Amer. J. Math. 60 (1938),889-896.
(4) LEVITAN, B.M.: Almost-periodic functions, Moscow, 1953 (in
Russian).
389-396.
(6) LEVITAN, B.M.: Mat. Zametki 1, no. 1 (1967), 35-44.
(7) GoRIN, E.A.: 'A function algebra invariant of the Bohr-van
Kampen theorem', Mat. Sb. 82 (124) (1970), 260-272 (in Rus-
sian). D. V Anosov
Editorial comments. Instead of the term 'average rota-
tion' one also uses 'mean motion', cf. e.g. [2], p. 310.
AMS 1980 Subject Classification: 42A75, 42A99
AVERAGE VALUE, THEOREM ON VARIATIONS
OF THE, in statistical mechanics - A statement on the
variation of the average value of a dynamic magnitude
in the Gibbs statistical aggregate as a result of an infini-
tesimal change in the Hamiltonian. The variation of the
average depends, generally speaking, on the way of
'inclusion' of the change in the Hamiltonian and on the
initial conditions. Let a system of many interacting par-
ticles (a quantum system or a classical system), that can
explicitly be described by a time-independent Hamil-
tonian H, be in thermodynamic equilibrium at the ini-
tial moment of time t~- 00. As a result of an adia-
batic inclusion of an infinitesimal, time-dependent per-
turbation of the Hamiltonian
H ~ H+8V(t),
where
8V(t) = ~e"-iE'8VE' >O, ~+O,
(E)
the Gibbs equilibrium average <A > of the explicitly
time-independent dynamic variable A will change (in a
linear approximation with respect to the perturbation)
by the magnitude
8<A(t = -2'7Ti~e't-lE' A 18VE~et),
(E)
 AVERAGING

where A 18VE~et) is the Fourier transform of 3) Multi-frequency non-autonomous systems

the retarded commutator Green function (cf. Green
function in statistical mechanics).
~ = p.X(x,y, t),
!!l!.. _
1 (3)
The main application of the theorem is in the theory dt - w(x,y, t)+p.Y(x,y, t).
of non-equilibrium irreversible processes (in which one
formulation of the theorem is also known as the Instead of the systems (1) - (3), 'simpler' averaged
fluctuation-dissipation theorem), and in deducing the systems of a first approximation are considered:
chains and systems of equations for the Green func-
tions from the chains and systems of equations for the ~; = p.Xo(x); (1 ')
correlation functions (cf. Correlation function in statisti-
~; = p.X1(X),}
cal mechanics).
(2')
References !!1.
dt
= w(x)'
[I) KUBO, R.: J. Phys. Soc. Japan 12 (1957),570.
[2) BOGOLYUBOV, JR., N.N. and SADOVNIKOV, B.I.: Zh. Eksper.
Teoret. Fiz., no. 43 (1962), 677. ~; = p.X2 (x, xo,Yo, t),}
[3) BOGOLYUBOV, JR., N.N. and SADOVNIKOV, B.I.: Some ques- (3')
tions in statistical mechanics, Moscow, 1975 (in Russian). !!1.-
dt - w(x, y, t),
[4) TYABLIKOV, S.V.: Methods of the quantum theory of magnetism,
Plenum Press, 1967 (translated from the Russian). where
1T
VN. Plechko f
Xo(x) = lim -T X (x, t, 0) dt,
T_oo 0
(4)
Editorial comments.
References (5)
[A1] BOGOLYUBOV, N.N. and BOGOLYUBOV, IR.N.N.,: Introduction 1 2" 2"
to quantum statistical mechanics, World Scientific, 1982. = (27T)" [ ... [X(X,Yl, ... ,Yn)dYI ... dYn,
AMS 1980 Subject Classification: 82A05
X 2 (x, xo,Yo, to) = (6)
to+T
AVERAGING - The operation of calculating the
mean values of functions forming part of the structure
= lim T1
T _00
f X(x, <j>(xo,Yo, to, t), t)dt.
10

of differential equations which describe periodic,

The formulas (4) - (6) express the most widespread
almost-periodic, and, generally, oscillating processes.
averaging methods.
The operation of averaging can be considered as a type
Formula (6) expresses the scheme of averaging 'along
of smoothing operator. Averaging methods first came
the generating solution'. In the function X(x,y, t), the
into use in celestial mechanics in the study of planetary
vector y is initially substituted by the generating solu-
motion around the Sun. Later, they propagated to a
tion of the system
wide variety of areas: to the theory of non-linear oscil- dx =0
lations, to physics, to the theory of automatic control, dt '
to astrodynarnics, and to others. Averaging methods !!l!..-
dt - w(x, Y, t),
have often provided approximate solutions for the ini-
tial equations. The most typical classes of differential after which the integral average (6) is calculated.
equations for which averaging methods are used are the The principal question which arises when the systems
following. (1) - (3) are changed is the construction of f-estimates
1) Standard systems in the sense of N.N. Bogolyubov for the norms
dx
dt = pX(x, t, p.), (1)
I x(t, p.)-x(t, p.) II, I y(t, p.)-y(t, p.) I
on the largest possible interval (of order 1/ fL) of time,
where x, X are vectors, t is the time, and fL is a small if
positive parameter. x (0, p.) = x(O, p.), Y (0, p.) = yeO, p.).
2) Multi-frequency autonomous 2'lT-periodic systems This is the essence of the problem of averagmg

~~ = p.X(X,y), 1 methods. For systems (1), this problem of averaging

methods was proposed by Bogolyubov, whose results
!
d
= w(x)+p.Y(x,y),
(2)
formed the basis of the modem algorithmic theory of
ordinary differential equations.
where x, y, X, Y are vectors, References
X(X,y+(27T)) X(x,y), Y(x,y+(2'lT)) Y(x,y), [I) BOGOLYUBOV, N.N. and MITROPOL'SKIi, YU.A.: Asymptotic
methods in the theory of non-linear oscillations, Gordon and
and w(x) is the frequency vector. Breach, 1961 (translated from the Russian).

313
AVERAGING
[2] MrrROPOL'SKIi, Yu.A.: Averaging methods in non-linear
mechanics, Kiev, 1971 (in Russian).
[3] VOLOSOV, V.M. and MORGUNOV, B.I.: Averaging methods in
the theory of non-linear oscillatory systems, Moscow, 1971 (in
Russian).
[4] GREBENIKOV, E.A. and RYABov, Yu.A.: Constructive methods
of analysis of non-linear systems, Moscow, 1979 (in Russian).
E.A. Grebenikov
Editorial comments. A recent up-to-date treatment of
averaging matters is given in [A 1].
References
[A1] SANDERS, J.A. and VERHULST, F.: Averaging methods in non-
linear dynamical systems, Springer, 1985.
AMS 1980 Subject Classification: 58FXX, 34C29
AXIAL VECfOR, pseudovector - A vector in oriented
space which, when the orientation of space is reversed,
is transformed into the inverse vector. An example of
an axial vector is the vector product of vectors.
BSE-3
AMS 1980 Subject Classification: 51 NXX, 15A72
AXIOM - A fundamental assumption, a self-evident
principle. In deductive scientific theories the axioms are
the fundamental primitive assumptions of the given
theory from which its remaining contents is derived by
deduction, i.e. by purely logical procedures. Cf.
Axiomatic method.
P. S. N ovikov
AMS 1980 Subject Classification: 03A05
AXIOM OF CHOICE - One of the axioms in set
theory. It states that for any family F of non-empty
sets there exists a function f such that, for any set S
from F, one has f(S)ES if is called a choice function
on F). For finite families F the axiom of choice can be
deduced from the other axioms of set theory (e.g. in the
system ZF).
The axiom of choice was explicitly formulated by E.
Zermelo (1904) and was objected to by many mathema-
ticians. This is e.\.plained, first, by its purely existential
character which makes it different from the remaining
axioms of set theory and, secondly, by some of its
implications which are 'unacceptable' or even contrad-
ict intuitive 'common sense'. Thus, the axiom of choice
implies: the existence of a Lebesgue non-measurable set
of real numbers; the existence of three subdivisions of
the sphere B
B = Uj U ... U Un,
B = Vj U ... U Vm ,
B = X U '" U Xn +m'
j AXIOM OF EXTENSIONALITY - One of the axioms
such that Vi is congruent with Xi, I~i~n, and Vj is
congruent with Xn +j, I ~j ~m. Thus, the sphere B is
divisible into a finite number of parts XI, ... , Xn +m
314
which can be moved in space to form two spheres
identical to it.
Many postulates equivalent to the axiom of choice
were subsequently discovered. Among these are: I) The
well-ordering theorem: On any set X there exists a total
order Rc;;,XXX such that any non-empty set Vex
contains a least element in the sense of the relation R;
2) the maximality principle (Zorn's lemma): If any
totally ordered subset V of a partially ordered set X is
bounded from above, X contains a maximal element; 3)
any non-trivial lattice with a unit element has a maxi-
mal ideal; 4) the product of compact topological spaces
is compact; and 5) any set X has the same cardinality
as XXx.
The axiom of choice does not contradict the other
axioms of set theory (e.g. the system ZF) and cannot be
logically deduced from them if they are non-
contradictory. The axiom of choice is extensively
employed in classical mathematics. Thus, it is used in
the following theorems. I) Each subgroup of a free
group is free; 2) the algebraic closure of an algebraic
field exists and is unique up to an isomorphism; and 3)
each vector space has a basis. It is also used in: 4) the
equivalence of the two definitions of continuity of a
function at a point (the t:-8-definition and the defini-
tion by limits of sequences) and in proving 5) the
countable additivity of the Lebesgue measure. The last
two theorems follow from the countable axiom of
choice (the formulation of the axiom includes the
condition of countability of the family F). It was
proved that the theorems I) to 5) are not deducible in
the system ZF if ZF is non-contradictory.
A model of set theory has been constructed which
meets the countable axiom of choice and in which each
set of numbers is Lebesgue-measurable. This model was
constructed on the assumption that the system ZF does
not contradict the axiom of the existence of an inacces-
sible cardinal number.
References
[I) FRAENKEL, A.A., BAR-HILLEL, Y. and LEVY, A.: Foundations of
set theory, North-Holland, 1973.
[2) JECH, TJ.: Lectures in set theory: with particular emphasis on
the method of forcing, Springer, 1971.
[3) JECH, TJ.: The axiom of choice, North-Holland, 1973.
v.N. Grishin
Editorial comments.
References
[A 1] JECH, TJ.: Settheory, Acad. Press, 1978.
AMS 1980 Subject Classification: 03E25, 04A25
of set theory, asserting that two sets are equal if they
contain the same elements:
\tu \tv (\tx (x EU=X EV) => U =v).

In a language not containing the equality symbol and
having only one predicate symbol E, the axiom of
extensionality has the form
Vu Vv (\fx (x EU~X EV) ~ Vz (u EZ~V EZ.
The axiom of extensionality has no real importace
for the formalization of mathematics in the
Zermelo- Fraenkel system ZF. Anything that can be
constructed within the system ZF can be formalized in
a system without the axiom of extensionality. Let ZF-
be the system obtained from ZF by removing the
axiom of extensionality and by replacing formulas of
the form u = v in the remaining axioms by the formula
Vx (XEU '> XEV).
Then it can be shown that there exists an interpretation
of ZF in ZF-. A similar assertion is valid for the
theory of types.
For Quine's system NF, obtained from the theory of
types by the 'erasure' of the type indices, the situation
is different: It is not possible to interpret NF in NP- .
The system NF- (NF without the axiom of exten-
sionality) is a rather weak system, and its consistency
can be proved in formal arithmetic. The sytem NF,
however, is not weaker than the theory of types with
the axiom of infinity.
References
[I] BARWISE,J. (ED.): Handbook of mathematical logic. North-
Holland, 1977.
[2] BOFFA, M.: 'The consistency problem for NF', J. Symbolic
Logic 42, no. 2 (1977),215-220.
VN. Grishin
Editorial comments.
References
[A1] ScOTI, D.S.: 'More on the axiom of extensionality', in Y.
Bar-Hillel, E.I.J. Poznanski and M.O. Rabin, et al. (eds.):
Essays on the foundation of mathematics, North-Holland,
1962.
AMS 1980 Subject Classification: 03E30, 03E65
AXIOM SCHEME - A unified manner of formulating
axioms (cf. Axiom) with the same syntactic structure. A
specific axiom scheme is usually realized by an expres-
sion ~ (which is most often formulated in a language
different from that in which the axioms are written),
which fixes the syntactic structure of the scheme, and
by using rules which make it possible to obtain any
axiom of the initial structure starting from ~.
In the context of previously formulated or unambigu-
ously understandable rules for generating axioms using
an expression ~, an axiom scheme is usually referred to
as the self-expression of ~. For instance, one speaks of
the axiom scheme a~(f3~a) of propositional calculus
P when one means the set of axioms of the form
A ~(B~A), where A and B are arbitrary formulas in
P.
AXIOMATIC METHOD
An example of a scheme of non-logical axioms is the
following variant of the induction scheme in the tradi-
tional axiomatizations of arithmetic:
(a(O)&V'x (a(x):,) a(x') :,) Vx a(x),
where a and x are assumed not to belong to the alpha-
bet of the language of the formalization considered and
are interpreted, respectively, as an arbitrary formula
and an arbitrary variable of this formalization.
The use of an axiom scheme usually permits one to
dispense with the rule of substitution when constructing
formal theories. Thus, in any sufficiently strong propo-
sitional calculus with two derivation rules - the rule of
substitution and the rule of conclusion - it is possible
to conduct proofs by substituting into axioms only,
which makes it possible to modify the calculus in an
equivalent manner by replacing each axiom by the
respective axiom scheme, and removing the rule of sub-
sitution as a derivation rule.
References
[I] KLEENE, S.c.: Introduction to metamathematics, North-Holland,
1951.
[2] CHURCH, A.: Introduction to mathematical logic, I, Princeton
Univ. Press, 1956.
F.A. Kabakov
AMS 1980 Subject Classification: 038XX, 03FXX
AXIOMATIC METHOD - A way of arriving at a
scientific theory in which certain primitive assumptions,
the so-called axioms (d. Axiom), are postulated as the
basis of the theory, while the remaining propositions of
the theory are obtained as logical consequences of these
axioms.
In mathematics, the axiomatic method originated in
the works of the ancient Greeks on geometry. The most
brilliant example of the application of the axiomatic
method - which remained unique up to the nineteenth
century - was the geometric system known as Euclid's
Elements (ca. 300 B.c.). At the time the problem of the
description of the logical tools employed to derive the
consequences of an axiom had not yet been posed, but
the Euclidean system was a very clear attempt to
obtain all the basic statements of geometry by pure
derivation based on a relatively small number of postu-
lates - axioms - whose truth was considered to be
self-evident.
The discovery of a non-Euclidean geometry by N.I.
LobachevskiI and J. Bolyai at the beginning of the
nineteenth century stimulated further development of
the axiomatic method. They showed that if the tradi-
tional, and apparently the only 'objectively true', fifth
postulate of Euclid concerning parallel lines is replaced
by its negation, then it is possible to develop in a
purely logical manner a geometry which is just as
elegant and meaningful as is Euclidean geometry. The
315
AXIOMATIC METHOD
attention of mathematicians of the nineteenth century
was thus drawn to the deductive manner of construct-
ing mathematical theories; this in turn gave rise to a
new problem, connected with the concept of the
axiomatic method itself and with the formal (axiomatic)
mathematical theory. With the gradually increasing
number of mathematical theories which had been
axiomatically derived - one can, in particular, mention
the axiomatic derivation of elementary geometry by M.
Pasch, G. Peano and D. Hilbert which, unlike Euclid's
Elements is logically unobjectionable, and Peano's first
attempt at the axiomatization of arithmetic - the con-
cept of a formal axiomatic system became more
rigorous (see below), resulting in a class of specific
problems which eventually established proof theory as
one of the main chapters of modern mathematical
logic.
It was recognized as early as the nineteenth century
that foundations must be created for mathematics and
for the relevant mathematical problems. It was mainly
in analysis that the basic concepts were rendered more
precise, and that more complex ideas were reduced to
simpler concepts by methods involving increasingly
rigorous logical reasoning (the (-() language of A.L.
Cauchy, the function-theoretic concepts of B. Bolzano
and K. Weierstrass, and the continuum of G. Cantor
and R. Dedekind). Also, the discovery of non-
Euclidean geometries stimulated the development of the
axiomatic method, the development of new ideas, and
the postulation of more general mathematical problems,
mainly those connected with concepts of an arbitrary
axiomatic theory, such as consistency, completeness
and independence of a given axiom system. The first
results in this direction were obtained by the method of
interpretation, which may be described as follows. Let
some specific mathematical object be assigned to each
primitive concept and to every relation of a given
axiomatic theory T. The population of such objects is
called the field of interpretation. It is now natural to
assign to each proposition IJr of T some sentence 1Jr'
about the elements of the field of interpretation, which
may be true or false. One then says that the proposi-
tion IJr of the theory T is true or false, respectively, in
this interpretation. The field of interpretation and its
properties are usually themselves the object of study of
some, usually different, mathematical theory T j, which
can also be axiomatic. The method of interpretation
can be used to establish relative consistency - i.e. to
demonstrate statements of the type 'if a theory T j is
consistent, so is the theory T' in the following manner.
Assume that the theory T has been interpreted in the
theory T j so that all axioms Ai of T are interpreted by
true statements A; of T j. All theorems of T, i.e. any
proposition A which is logically deducible from the
316
axioms Aj in T, is interpreted in T 1 by some proposi-
tion A, which is deducible in T 1 from the interpreta-
tions A; of the axioms Aj , and is therefore true. This
last statement is based on an additional implicit
assumption to the effect that the logical tools of the
theories T and T 1 are to some extent similar, in prac-
tice this condition is usually satisfied. (In the initial
applications of the method of interpretation, this
assumption was not even discussed, since it was
assumed to be self-evident; in fact, during the first
attempts at demonstrating a relative consistency
theorem, the logical tools of T and T 1 were simply
identical; this was the classical theory of predicates.)
Now let the theory T be inconsistent, i.e. let both some
statement A of this theory and its negation be derivable
in it. It then follows from the above that the proposi-
tions A and non-A will both be true statements of
T 1, i.e. T 1 is inconsistent. This was the method
employed by F. Klein and H. Poincare to show that
LobachevskiI's non-Euclidean geometry is consistent if
Euclid's geometry is consistent; also, the problem of the
consistency of Hilbert's axiomatization of Euclidean
geometry was reduced by Hilbert to the problem of the
consistency of arithmetic. The method of interpretation
can also be used to solve the problem of independency
of a systems of axioms: In order to prove that an
axiom A of a theory T is independent of the other
axioms of that theory, i.e. cannot be deduced from
them, and that its inclusion is therefore essential in
order to obtain the complete theory, it is sufficient to
construct an interpretation of T in which A is false,
while all the other axioms are true. Another form of
this manner of proving the independence of A is to
establish the consistency of the theory obtained by
replacing A in T by its negation. The result of the
above-mentioned reduction of the problem of the con-
sistency of Lobachevskll's geometry to the problem of
the consistency of Euclidean geometry, and of reducing
the problem of the consistency of Euclidean geometry
to the problem of the consistency of arithmetic, is the
conclusion that Euclid's fifth postulate is not derivable
from the remaining geometrical axioms, provided that
the arithmetic of positive integers is consistent. The
weak point of the method of interpretation is the fact
that, as far as problems of consistency and indepen-
dence of axiom systems are concerned, it inevitably
yields results of merely a relative character. However,
its important achievement is the fact that it served to
reveal that arithmetic plays a special role, in that it is a
mathematical theory to whose consistency the problem
of consistency of several other theories is reduced.
The method of the so-called formalism of the founda-
tions of mathematics, due to Hilbert and his school,
was a further step (and, in a sense, a peak) in the

development of the method. It rendered the concept of
an axiomatic theory more precise by introducing the
notion of a fonnal system as the next stage in the
development of the axiomatic method. As a result, it
proved possible to treat the mathematical theories
themselves as precise mathematical objects and to con-
struct a general theory of such theories - a so-called
meta-theory. The possibility of solving all the main
problems in the foundations of mathematics in this way
appeared very attractive, and Hilbert himself was
tempted to follow this path. The principal concept in
this approach was that of a formal system. Any formal
system is constructed as a perfectly definite class of
expressions - that is, formulas from which a subclass
of formulas, known as the theorems of the formal sys-
tem, is derived in a certain well-defined manner. The
formulas of the formal system are not directly meaning-
ful, and in general they may be constructed from signs
or symbols merely chosen for reasons of technical con-
venience. In fact, the way of constructing formulas and
the concept of a theorem of some formal system are so
chosen that the resulting formal apparatus can be used
for the most adequate expression of some mathematical
(or even non-mathematical) theory, more precisely, as
an expression of both its factual contents and its
deductive structure. The general procedure for con-
structing an arbitrary formal system S is as follows.
I. The language of the system S.
a) The alphabet is a list of elementary symbols of the
system.
b) The construction (syntactic) rules are rules accord-
ing to which the formulas of S are constructed from the
elementary symbols. A sequence of elementary sym-
bols is a taken as a formula if and only if it can be
constructed by means of syntactic rules.
II. The axioms of the system S. A certain set of for-
mulas (usually a finite or countable set), said to be the
axioms of the system S, is specified.
III. The derivation rules of the system S (or deduction
rules). A (usually finite) population of predicates
R;, ... ,Rk is defined on the set of all formulas of S.
Let Rj(x I, . . . ,xn, + I ) be one of these predicates
(nj>O); if, for given formulas F I , ... ,Fn,+I, the
statement R;(FI' ... ,Fn,+d is true, then one says that
the formula Fn, + I follows directly from the formulas
F 1> ,Pn, by the rule R;.
Specifying I, II and III is the same as specifying the
formal system S as a precise mathematical object, since
the concept of a theorem or a derivable formula of S is
formed for all formal systems in a uniform manner (the
degree of accuracy depending on the level of accuracy
of the alphabet, of the syntactic rules and of the deriva-
tion laws, i.e. of the predicates R 1> ,Rk ). A deriva-
AXIOMATIC METHOD
tion in S is any finite sequence of formulas in that sys-
tem in which each formula is either an axiom of S or
follows directly from some formulas which precede it in
the sequence in accordance with one of the derivation
laws R; of S. A formula of S is called a theorem of this
system if there exists a derivation in S ending in this
formula.
Any specific mathematical theory T may be
translated into the language of a suitable formal system
S so that any meaningful (true or false) statement of T
can be expressed by one of the formulas of S.
It can naturally be expected that this method of for-
malization would make it possible to construct all the
meaningful elements of any mathematical theory on the
precise and apparently reliable basis represented by the
concept of a derivable formula (a theorem of the formal
system), and to solve fundamental problems such as the
problem of the consistency of the mathematical theory
by proving the corresponding statements about the for-
mal system that formalizes this theory. Since the formal
systems of the type just described are exact, or 'finitis-
tic' , using the term used by the school of Hilbert,
mathematical objects, it could be expected that it would
be possible to obtain jinitistic proofs of consistency
statements, i.e. proofs which would in a certain sense
be effective, that is, independent of such powerful tools
as, for example, the abstraction of actual infmity (which
is one of the reasons for the difficulties encountered in
the foundations of classical mathematical theories).
Thus, the requirement that the tools employed in arriv-
ing at results concerning formal systems, and in partic-
ular concerning their consistency, be finitistic,
represented a regular feature of Hilbert's formalistic
program. However, the results obtained by K. GOdel in
the early nineteen-thirties frustrated the main expecta-
tions based on this program. Godel showed that:
. 1) Any natural consistent formalization of arithmetic
or of any other mathematical theory which involves
arithmetic (e.g. set theory) is incomplete and cannot be
completed in the sense that a) S contains (true with
respect to their contents) unsolvable formulas, i.e. for-
mulas A such that neithtr A nor non-A are derivable in
S (the incompleteness of formal arithmetic); b) whatever
the finite set of supplementary axioms (e.g. of unsolv-
able formulas in S) which are introduced to enlarge S,
the new, stronger, system thus produced inevitably con-
tains its own unsolvable formulas (incompleteness; see
[5] and GOdeI incompleteness theorem).
2) If formal arithmetic is in fact consistent then,
while the statement of its consistency is expressible in
its own language, it is impossible to prove this state-
ment by the methods of formalized arithmetic itself.
This means that, even in the case of arithmetic, it is
intrinsically impossible to exhaust all its contents-wise
317
AXIOMATIC METHOD
true statements by a class of derivable formulas of any
formal system, and that there is no hope of obtaining
any finitistic proof of the consistency of arithmetic,
since it seems that any reasonable precision of the con-
cept of a finitistic proof is formalizable in formal arith-
metic.
All this imposes definite limitations on the possibili-
ties of the axiomatic method in the form it takes in the
framework of Hilbert's formalism. However, within
these limits it continues to play an important role in
the foundations of mathematics. Thus, subsequent to
the publication of GOdel's results, he himself (1938 -
1940) and P. Cohen (1963), basing themselves on the
axiomatic approach with the use of the method of
interpretation, obtained fundamental results concerning
the compatibility and mutual independence of the
axiom of choice and the continuum hypothesis in set
theory [6], [7]. As regards consistency, which is a funda-
mental problem in the foundations of mathematics, it
had become clear in view of GOdel's results that it can
be solved only with the aid of tools and ideas which are
not finitistic. Various approaches to this question
proved possible, not all of which are equally acceptable
to mathematicians, mainly because of the differences of
opinion as to the acceptability of various logical tools.
Among the results obtained concerning the consistency
of formal systems, one can mention, first of all, the
proof of the consistency of formal arithmetic [8], based
on infinite induction up to a certain countable transfin-
ite number. Another, more recent, example of this kind
is the attempt to use certain ideas of intuitionism for
proving the consistency of the formal system of analysis
[9].
References
[1] HEATH, T.L.: The thirteen books Euclid's elements, 11-13,
Dover, reprint, 1956. (Translated from the Greek).
[2] KAGAN, V.F.: The foundations of geometry, 1, Moscow-
Leningrad, 1949 (in Russian).
[3] HILBERT, D.: Grundlagen der Geometrie, Springer, 1913.
[4] llEANo, G.: 'Principii di logica matematica', Rivista Matema-
tika 1 (1891), 1-10.
[5] GoDEL, K.: 'Ueber formal unentscheidbare Satze der Principia
Mathematica und verwandter Systeme 1', Monatsh. Math. Phy-
sik 38 (1931), 178-198.
[6] GoDEL, K.: The consistency of the axiom of choice and of the
generalized continuum hypothesis with the axioms of set theory,
Princeton Univ. Press, 1940.
[7] CoHEN, PJ.: Set theory and the continuum hypothesis, Benja-
min, 1966.
[8] GENTZEN, G.: 'Die Widerspruchsfreiheit der reinen
Zahlentheorie', Math. Ann. 112 (1936),493-565.
[9] SPECTOR, C.: 'Provable recursive functionals of analysis', in
Recursive function theory, Amer. Math. Soc., 1962, pp. 1-27.
P.S. Novikov
Editorial comments.
References
[A1] KLEENE, S.C.: Introduction to metamathematics, North-
Holland & Noordhoff, 1959.
AMS 1980 Subject Classification: 03A05, 00A25
318
AxIOMATIC SEI' 1HEORY - The branch of
mathematical logic in which one deals with fragments
of the informal theory of sets by methods of mathemat-
ical logic. Usually, to this end, these fragments of set
theory are formulated as a formal axiomatic theory. In
a more narrow sense, the term 'axiomatic set theory'
may denote some axiomatic theory aiming at the con-
struction of some fragment of informal ('naive') set
theory.
Set theory, which was formulated around 1900, had
to deal with several paradoxes from its very beginning.
The discovery of the fundamental paradoxes of G. Can-
tor and B. Russell (cf. Antinomy) gave rise to a
widespread discussion and brought about a fundamen-
tal revision of the foundations of mathematical logic.
The axiomatic direction of set theory may be regarded
as an instrument for a more thorough study of the
resulting situation.
The construction of a formal axiomatic theory of sets
begins with an accurate description of the language in
which the propositions are formulated. The next step is
to express the principles of 'naive' set theory in this
language, in the form of axioms and axiom schemes. A
brief description of the most widespread systems of
axiomatic set theory is given below. In this context, an
important part is played by the language which con-
tains the following primitive symbols: 1) the variables
x,y, z, U, V, Xl> , which play the part of common
names for the sets in the language; 2) the predicate
symbols E (sign of incidence) and = (sign of equality);
3) the description operator L, which means 'an object
such that ... '; 4) the logical connectives and quantif-
iers: B (equivalent), -+ (implies), V (or), /\ (and), -,
(not), 'fI (for all), 3 (there exists); and 5) the
parentheses ( and ). The expressions of a language are
grouped into terms and formulas. The terms are the
names of the sets, while the formulas express proposi-
tions. Terms and formulas are generated in accordance
with the following rules.
RI. If T and 0 are variables or terms, then (TEO) and
(T=O) are formulas.
Rl. If A and B are formulas and x is a variable, then
(A BB), (A -+B), (A V B), (A I\B), -,A, 'fI xA, 3 xA are
formulas and LXA is a term; the variable x is a term.
For instance, the formula 'fIX(XEy-+XEZ) is tan-
tamount to the statement y is a subset of z', and can
be written as y C;:z; the term LW'fIy(y EWBY c;:z) is the
name of the set of all subsets z and, expressed in con-
ventional mathematical symbols, this is Pz. Let the
symbol <9 mean 'the left-hand side is a notation for the
right-hand side'. Below a number of additional nota-
tions for formulas and terms will be presented.
The empty set:
o # u'Vy...,yex.
The set of all x such that A (x):
 AXIOMATIC SET THEORY

{x: A(x)} ~ LZ'VX(XEZ _A(x, axiomatic system of this type was the system Z, due to
where z does not enter freely in A (x) (i.e. IS not a E. Zermelo (1908). However, this system does not allow
parameter of the formula A (x. a natural formalization of certain branches of
The unordered pair x and y: mathematics, and the supplementation of Z by a new
principle - the axiom of replacement - was proposed
{x,y} ~ {z: Z =x V z=y}. by A. Fraenkel in 1922. The resulting system is known
The single-element set consisting of x: as the Zermelo- Fraenkel system and is denoted by ZF.
{x} ~ {x,x}. b) The second group is constituted by systems the
axioms of which are selected in the context of giving
The ordered pair x and y:
some explanations for paradoxes, for example, as a
<x,y> ~ {{x}, {x,y}}. consequence of non-predicative definitions. The group
The union of x and y: includes Russell's ramified theory of types, the simple
theory of T-types, and the theory of types with
x Uy ~ {z: ZEX V zEy}.
transfinite indices (cf. Types, theory of).
The intersection of x and y: c) The third group is characterized by the use of
x ny ~ {z: ZEX I\ZEy}. non-standard means of logical deduction, multi-valued
The union of all elements of x: logic, complementary conditions of proofs and infinite
derivation laws. Systems in this group have been
Ux ~ {z: 3V(ZEV VVEX)}.
developed to the least extent.
The Cartesian product of x and y: d) The fourth group includes modifications of sys-
x Xy ~ {z: 3uv(z = <u, v> 1\ UEX 1\ v Ey)}.
tems belonging to the first three groups and is aimed at
attaining certain logical and mathematical objectives.
Notation for: w is a function: Only the system NBG of Neumann - Gooel- Bemays
Fnc(w) ~ 3v(w C;v Xv) 1\ (1925) and the system NF of W. Quine (1937) will be
1\ 'VUVI V2U, VI > EW 1\ <U, V2>EW ~ VI =V2)' mentioned here. The construction of the system NBG
was motivated by the desire to have a finite number of
The values of the function w on the element x:
axioms of set theory, based on the system ZF. The sys-
W'X ~ o/<X,y>EW. tem NF represents an attempt to overcome the stratifi-
The standard infinite set z cation of the concepts in the theory of types.
The systems Z, ZF and NF can be formulated in the
Inf(z) ~ 0 EZ 1\ 'VU(UEZ ~u U {U}EZ).
language described above. The derivation rules, and
The axiomatic theory A that follows is the most com- also the so-called logical axioms, of these systems are
plete representation of the principles of 'naive' set identical, and form an applied predicate calculus of the
theory. The axioms of A are: first order with equality and with a description opera-
AI. Axiom of extensionality: tor. Here are the axioms of equality and of the descrip-
'VX(X EY_X EZ) ~ y=z tion operator:
('if the sets x and y contain the same elements, they are x = x, x =y ~ (A(x)~A(y,
equal'); where A (x) is a formula not containing the bound vari-
A2. Axiom scheme of comprehension: able y (i.e. it has no constituents of the type '1y, 3y, ry),
3y 'Vx(x E y_A), while A (y) is obtained from the formula A (x) by
replacing certain free entries of the variable x with y:
where A is an arbitrary formula not containing y as a
parameter ('there exists a set y containing only elements 3!xA(x) ~ A (tXA(x,
x for which A'). where the quantifier 3!x means that 'there exists one
This system is self-contradictory. If, in A2, the for- and one only x', while the formula A(tXA(x is
mula -,x EX IS taken as A, the formula obtained from the formula A (x) by replacing all free
'1X(XEJ~XEX) readily yields yEy~Ey, which is
entries of the variable x with the term tXA (x). The
a contradiction. quantifier 3!x can be expressed in terms of the quantif-
The axiomatic systems of set theory may be subdi- iers 'v' and 3 and equality.
vided into the following four groups. Non-logical axioms of the system Z:
a) The construction of axiomatic systems in the first ZI.The axiom of extensionality AI.
group is intended to restrict the comprehension axioms Z2. The pair axiom:
so as to obtain the most natural means of formalization
3U'VZ(ZEU _ z=x V z=y)
of conventional mathematical proofs and, at the same
time, to avoid the familiar paradoxes. The first ('the set {x, y} exists');

319
AXIOMATIC SET THEORY

Z3. The union axiom:
3yV'X(XEy ~ 3t(tEz I\XEt
('the set U z exists');
1.4. The power set axiom:
3y V'x(xEJ ~ x ~z)
('the set pz exists');
Z5. The separation axiom scheme:
3yV'x(xEJ ~ XEZ 1\ A (x
('there exists a subset Z consisting of the elements x in z
for which A (x) is true'); the axioms Z2 - Z5 are exam-
ples of axioms of comprehension;
Z6. The axiom of infinity:
3z Inf(z);
Z7. The axiom of choice:
V'z 3w(Fnc(w) 1\ V'x(x EZ 1\ oX = 0 ~ w'x EX
('for any set z there exists a function w which selects,
out of each non-empty element x of the set z, a unique
element w'x'). The above axioms are complemented by
the regularity axiom:
ZS.
V'X(oX =0 ~ 3y(yEX I\y n X = 0,
which is intended to postulate that there are no des-
cending chains x2EX),x3EX2,X4EX3, .... Axiom ZS
simplifies constructions in Z, and its introduction does
not result in contradictions.
The system Z is suitable for developing arithmetic,
analysis, functional analysis and for studying cardinal
numbers smaller than ~w. However, if the alephs are
defined in the usual manner, it is no longer possible to
demonstrate the existence in Z of ~w and higher cardi-
nal numbers.
The system ZF is obtained from Z by adding
Fraenkel's replacement axiom scheme, which may be
given in the form of the comprehension axiom scheme:
ZF9.
3y \tx(x Ey ~ 3v(v EZ 1\ x =LtA (t, v)
('there exists a set y consisting of x, x = LtA (I, v), where
v runs through all the elements of a set z'). In other
words, y is obtained from z if each element v of z is
replaced with LtA(/, v).
The system ZF is a very strong theory. All ordinary
mathematical theorems can be formalized in terms of
ZF.
The system NBG is obtained from ZF by adding a
new type of variables the class variables
X, Y, Z, . .. - and a finite number of axioms for
forming classes, by means of which it is possible to
prove formulas of the type
3y\tX(XEY~A(x,
where A (x) is a formula of NBG which does not coun-
tain bound class variables or the symbol t. Since any
formula A (x) can be used to form a class, the infinite
number of ZFaxioms can be replaced by a finite
number of axioms containing a class variable. The
axiom of choice has the form:
3X(Fnc(X) 1\ \tx(-,x = 0 ~ X' X EX
and confirms the existence of a selection function,
which is unique for all sets and which constitutes a
class.
The system NF has a simpler axiomatic form, viz.: 1)
the axiom of extensionality; and 2) the axioms of
comprehension in which a formula A can be stratified,
i.e. it is possible to assign to all variables of the for-
mula A superscript indices so as to obtain a formula of
the theory of T-types, i.e. in the subformulas of type
x E y the index of x is one lower than the index of y.
The system NF has the following characteristics:
a) the axiom of choice and the generalized continuum
hypothesis are disprovable;
b) the axiom of infinity is demonstrable (d. Infinity,
axiom of);
c) the extensionality axiom plays a very important
role. Thus, if the extensionality axiom is replaced by
the slightly weaker axiom:
(3u(uEJ) 1\ \tU(UEJ~UEZ ~ y=z,
which permits a large number of empty sets, while the
comprehension axioms of NF remain unchanged, a
fairly weak theory is obtained: The consistency of the
resulting system can be proved even in formal arith-
metic.
Results concerning the interrelationships between the
systems just described are given below.
a) Any formula of ZF is demonstrable in NBG if
and only if it is demonstrable in ZF.
f3) In ZF it is possible to establish the consistency of
Z, completed by any finite number of examples of the
axiom scheme of replacement ZF9. Thus, ZF is much
stronger than Z.
y) The consistency of T is demonstrable in Z, so that
Z is stronger than T.
{) NF is not weaker than T in the sense that it is
possible to develop the entire theory of types in NF.
The axiomatic approach to the theory of sets has
made it possible to state a proposition on the unsolva-
bility in principal (in an exact sense) of certain
mathematical problems and has made it possible to
demonstrate it rigorously. The general procedure for
the utilization of the axiomatic method is as follows.
Consider a formal axiomatic system S of the theory of
sets (as a rule, this is ZF or one of its modifications)
that is sufficiently universal to contain all the conven-
tional proofs of classical mathematics, and for all ordi-
nary mathematical facts to be deducible from it. A

320

given problem A may be written down as a formula in
the language S. It is then established by mathematical
methods that neither A nor its negation can be deduced
in S. It follows that problem A cannot be solved (in
either way) by tools of the theory S, but since this
theory S was assumed to contain all ordinary methods
of proof, the result means that A cannot be solved by
ordinary methods of conclusion, i.e. A is 'transcenden-
tal'.
Results which state that a proof cannot be performed
in the theory S are usually obtained under the assump-
tion that S, or some natural extension of S, is con-
sistent. This is because on the one hand, the problem
can be non-deducible in S only if S is consistent, but
such consistency cannot be established by the tools
offered by S (cf. GOOei incompleteness theorem), i.e.
cannot be derived by ordinary methods. On the other
hand, the consistency of S is usually a very likely
hypothesis; the very theory S is based on its truth.
Furthermore, the axiomatic approach to the theory
of sets made it possible to accurately pose and solve
problems connected with effectiveness in the theory of
sets, which had been intensively studied during the ini-
tial development of the theory by R. Baire, E. Borel, H.
Lebesgue, S.N. Bernstein [S.N. Bernshtein], N.N. Luzin
and W. Sierpiilski. It is said that an object in the
theory of sets which satisfies a property W is effectively
defined in the axiomatic theory S if it is possible to
construct a formula A (x) of S for which it can be
demonstrated in S that it is fulfilled for a unique
object, and that this object satisfies property W.
Because of this definition it is possible to show in a
rigorous manner that for certain properties W in S it is
impossible to effectively specify an object which satis-
fies W, while the existence of these objects in S can be
established. But since the chosen theory S is sufficiently
universal, the fact that the existence of certain objects
in S is ineffective is also a proof of the fact that their
existence cannot be effectively established by ordinary
mathematical methods.
Finally, the methods of the axiomatic theory of sets
make it possible to solve a number of difficult prob-
lems in classical branches of mathematics as well: in
the theory of cardinal and ordinal numbers, in descrip-
tive set theory and in topology.
Some of the results obtained by the axiomatic theory
of sets are given below. Most of the theorems concern
the axiomatic set theory of Zermelo- Fraenkel (ZF),
which is now the most frequently employed. Let ZF-
be the system ZF without the axiom of choice Z7. In
view of a), the results can be readily adapted to the
system NBG as well.
1) It was shown in 1939 by K. GOdel that if ZF- is
consistent, it will remain consistent after the axiom of
AXIOMATIC SET THEORY
choice and the continuum hypothesis have been added.
It follows that it is impossible to disprove the axiom of
choice or the continuum hypothesis in ZF. In order to
prove this result, Godel constructed a model of the
theory ZF consisting of the so-called GOdel construc-
tive sets (cf. GOOel constructive set), this model plays an
important role in modern axiomatic set theory.
2) The problem as to whether or not the axiom of
choice or the continuum hypothesis is deducible in ZF
remained open until 1963, when it was shown by P.J.
Cohen, using his forcing method, that if ZF- is con-
sistent, it will remain consistent after the addition of
any combination of the axiom of choice, the continuum
hypothesis or their negations. Thus, these two problems
are independent in ZF.
The principal method used for establishing that a
formula A is not deducible in ZF is to construct a
model of ZF containing the negation of A. Cohen's
forcing method, which was subsequently improved by
other workers, strongly extended the possibilities of
constructing models of set theory, and now forms the
basis of almost all subsequent results concerning non-
deducibility. For instance:
3) It has been shown that one can add to ZF,
without obtaining (additional) inconsistencies, the
hypothesis stating that the cardinality of the set of sub-
sets of a set x may be an almost arbitrary pre-given
function of the cardinality of x on regular cardinals
(the only substantial restrictions are connected with
Konig's theorem).
4) M.Y a. Suslin (1920) formulated the following
hypothesis. Any linearly totally ordered set such that
any pairwise non-intersecting family of non-empty open
intervals in it is 'at most countable must contain a
countable everywhere-dense subset. The non-
deducibility of Suslin's hypothesis in ZF was esta-
blished by Cohen's method.
5) It was shown that the following postulate: 'Any
subset of the set of real numbers is Lebesgue measur-
able' is unsolvable in ZF- (without the axiom of
choice).
6) The interrelationship of many important problems
of descriptive set theory with ZF was clarified. The first
results relating to this problem were demonstrated by
P.S. Novikov [5]. The methods of axiomatic set theory
made it possible to discover previously unknown con-
nections between the problems of 'naive' set theory. It
was proved, for example, that the existence of a Lebes-
gue non-measurable set of real numbers of the type ~1
(i.e. A 2) implies the existence of an uncountable III
(i.e. C s1) set without a perfect subset.
7) It was proved that an effectively totally ordered
continuum is absent in ZF. Numerous results proved
the absence of effectively defined objects in the descrip-
321
AXIOMATIC SET THEORY
tive theory of sets and m the theory of ordinal
numbers.
References
[1) FIlAENKEL, A.A., BAR-HILLEL, Y. and LEVY, A.: Fowuiations of
set theory, North-Holland, 1973.
(2) COHEN, P.J.: Set theory and the continuum hypothesis, Benja-
min, 1966.
(3) JEeR, T.J.: Lectures in set theory: with particular emphasis on
the method offorcing, Springer, 1971.
(4) DRAKE, F.R.: Set theory: an introduction to large cardinals,
North-Holland, 1974.
(5) NOVIKOV, P.S.: 'On the consistency of certain propositions of
the descriptive theory of sets', Amer. Math. Soc. Transl. 29
(1963), 51-89. (Trudy Mat. Inst. Steklov. 38 (1951),279-316)
V.N. Grishin
A. G. Dragalin
Editorial comments. Gbdel's book [A4] contains his
proof of the statement in 1) above.
References
[Al] JOCR, T.J.: Set theory, Acad. Press, 1978.
[A2] lEMMON, EJ.: Introduction to axiomatic set theory, Routledge
& Kegan Paul, 1968.
[A3] TAKEUTI, G. and ZARING, W.M.: Introduction to axiomatic set
theory, Springer, 1971.
[A4] GOOEL, K.: The consistency of the continuum hypothesis,
Princeton Univ. Press. 1940.
AMS 1980 Subject Classification: 03EXX, 04-XX
AxIOMATIZED CLASS - A class of models of one
type, defined by an axiom system. A class K of models
of a formal language L is said to be axiomatized (fin-
itely axiomatized) if there exists a (finite) system ~ of
closed formulas of L such that K contains those and
only those models on which all formulas of ~ are
defined and are true (cf. Algebraic system). A class of
models of a recursive signature is said to be recurSively
axiomatized if it can be specified by a recursive set of
axioms.
Many classes of algebraic systems studied in
mathematics are defined by a system of axioms of a
first-order language. For instance, the cl~)ses of all
Boolean algebras, all groups, all fields, and all lattices
are finitely axiomatized. The classes of all torsion-free
groups, all fields of characteristic zero and all algebrai-
cally closed fields are recursively axiomatized, but not
necessarily finitely axiomatized. The theory of axioma-
tized classes reveals regularities common to all classes
of objects defined by a specific language; it has been
well developed for first-order languages, and therefore
only such classes and formulas will be dealt with in
what follows.
Two models are said to be elementarily equivalent if
any formula of the first-order language which is true in
one of them is also true in the other. A model we is
said to be an elementary extension of a model ~ if any
formula which is defined and is true in ~ is also true in
Wl.
322
An elementary closed class K of models is called
complete if all its models are elementarily equivalent.
Every axiomatized class of models is a sum of
pairwise-disjoint complete classes. A class is said to be
categorical in cardinality m if all its models of cardinal-
ity m are isomorphic. A complete class of models of a
countable signature that is categorical in an uncount-
able cardinality, is categorical in all uncountable cardi-
nalities, but may be non-categorical in a countable car-
dinality; in such a case the class has a countable
number of pairwise non-isomorphic countable models.
For any n=l=2 there exists a complete axiomatized class
with exactly n non-isomorphic countable models.
An axiomatized class K of models is called solvable if
there exists an algorithm by which it is possible to tell,
for any closed formula of the language L, if it is true or
false for each model in K. The following theorem
describes the interconnection of complete, categorical
and solvable classes: If K is categorical of infinite car-
dinality and has no finite model, then it is complete. A
recursively-complete axiomatized class of models is
solvable.
Reductive and projective classes are generalizations
ofaxiomatized clases. Projective classes are defined by a
second-order axiom in the form
3T1 3Tn 2!(PJ, ... ,Ph Tlo ... ,Tn),
where Pi, Tj are predicate variables, and
2l(Pj, ... ,Pk , Tj, ... ,Tn) is a formula of the signa-
ture o=(Pj, ... ,Pk> Tj, ... ,Tn). Many properties of
axiomatized classes can be applied to these classes.
References
[11 MAL'TSEV, A.I.: Algebraic systems, Springer, 1973 (translated
from the Russian).
[2] MAL'TSEV, A.I.: Proc. 4-th All-Union Math. Congress 1961, Vol.
1, Leningrad, 1963, pp. 169-198 (in Russian).
(3) COHN, P.M.: Universal algebra, Reidel, 1981.
A.D. Taimanov
Editorial comments. An axiomatized class is also called
an axiomatic class. Elementary equivalent is also called
indiscernible.
AMS 1980 Subject Classification: 03CXX
AxONOMETRY - One of the methods of mapping
three-dimensional figures into a plane. An axonometry
consists in projecting the figure, after an orthogonal
Cartesian coordinate system and a coordinate plane
into which the figure will be projected have been
chosen, on the chosen plane for the diagram (figure).
Depending on the mode of the projection, the axo-
nometry may be parallel or central. If the direction of
the parallel projection into the chosen plane is perpen-
dicular to this plane, the axonometry is said to be nor-
mal or orthogonal; otherwise it is said to be skew. The
parameters of an orthogonal axonometry are the
 AXONOMETRY

cosines of the angles between the coordinate axes and theorem: Any tetrahedron can be projected to an arbi-
the chosen plane, and are sometimes called distortion trarily given plane tetragon.
indices. If two distortion indices are equal, the axo-
References
nometry is called a dimetry, and if three distortion [I] GLAZUNov, E.A. and CmrrvERUKHIN, N.F.: Axonometry, Mos-
indices are equal, an isometry; if all distortion indices cow, 1953 (in Russian). AB I
are different, it is called a trimetry. The principal .. vanov
theorem on axonometries is the Pohlke- Schwartz AMS 1980 Subject Classification: 51 M15

323
 B
/l-DISTRIBUTION - See Beta-distribution. BAER MULTIPLICATION - A binary operation on
AMS 1980 Subject Classification: 62EXX the set of classes of extensions of modules, proposed by
R. Baer [1]. Let A and B be arbitrary modules. An
(B, <p)-STRUCTURE - A structure on a vector bun- extension of A with kernel B is an exact sequence:
dle (or sphere bundle, etc.) that is a generalization of o --4 B --4 X --4 A --4 O. (I)
the concept of the structure group of a fibration.
The extension (1) is called equivalent to the extension
Let CPn: Bn~BOn be a fibration and let ~ be an n-
dimensional vector bundle over a space X, classified by 0--4 B --4 XI --4 A --40

the mapping ~: X ~BOn' Then the homotopy class lift- if there exists a homomorphism a: X~Xl forming part
ing the mapping ~: X ~BOn to a mapping in Bn is of the commutative diagram
called a (Bn' cpn)-structure on ~, i.e. it is an equivalence
class of mappings r: X ~Bn such that CPn or =~, where
two mappings rand r': X~Bn are said to be
equivalent if they are fibrewise homotopic. No method
of consistently defining (Bn' <pn)-structures for
equivalent fibrations exists, because this consistency
depends on the choice of the equivalence. The set of equivalence classes of extensions is denoted
Let there be a sequence (B, cP, g) of fibrations by Ext(A, B). The Baer multiplication on Ext(A, B) is
CPr: Br~BOr and mappings gr: Br~Br+l such that induced by the operation of products of extensions
jrocpr=CPr+logr (jr:BOr~BOr+l is the standard map- defined as follows. Let
ping). The family {B" cP" gr} (and sometimes only f3 a

(B" CPr is called a structure series. An equivalence class o --4 B --4 X --4 A --40, (2)
of sequences of (B" CPr )-structures on the normal bundle
(3)
{~r} of a manifold M n is called a (B, cP )-structure on M;
they coincide beginning from some sufficiently large r. be two extensions. In the direct sum XEB Y the submo-
A manifold M n with a fixed (B, cp)-structure on it is dules
c= {(x,y): a(x)=al(Y)}
called a (B, cp)-manifold.
Instead of BOn, a more general space BGn, classify- and
D = {(-x,y): x=/3(b),y=/3I(b)}
ing sphere bundles, can be considered and (B, cp)-
are selected. Clearly, DeC, so that one can define the
structures can be introduced on them.
quotient module Z = C / D. The Baer product of the
References extensions (2) and (3) is the extension
[1] LASHOF, R.: 'Poincare duality and cobordism', Trans. Amer.
Math. Soc. 109 (1963), 257-277.
[2] STONG, R.E.: Notes on cobordism theory, Princeton Univ. Press,
1968. where
Yu.B. Rudyak /32(b) = [/3(b), O] = [O,/3'(b)],
Editorial comments. Here and
BOn = lim infGrasn(R'+n)
r~oo

References
is the limit of the Grassmann manifolds of r-planes in R n + r [I] BAER, R.: 'Erweiterung von Gruppen und ihren lsomorphis-
AMS 1980 Subject Classification: 55R25, 55R35 men', Math. Z. 38 (1934).374-416.

325
BAER MULTIPLICATION
[2) CARTAN, H. and ElLENBERG, S.: Homological algebra, Prince-
ton Univ. Press, 1956.
VE. Govorov
AMS 1980 Subject Classification: 18G 15
BAlKE CLASSES - Families of real functions which
are defined inductively using the ordinal number of
limit transitions involved in the definition of the func-
tion, and which constitute the classification of functions
proposed in 1899 by R. Baire [1], known as the Baire
classification. The set of all continuous functions
f : A ~R, where A is a metric space, is called the zero-
th Baire class H o. The first Baire class HI is the set of
discontinuous functions f: A ~R which are the limit of
a sequence of continuous functions converging at each
point. The Baire class H a' where a is an ordinal
number of the first or second class, is defined as the set
of functions f: A ~R which do not form part of any of
the preceding classes, but which may be represented as
f=limn~ooJ,., where J,. EH P.' f3n<a. The union of the
Baire classes H a over all ordinal numbers of the first
and second class constitutes the set of Baire functions.
This is the minimal class of functions f :A ~R contain-
ing all the continuous functions and closed with respect
to pointwise convergence. A linear combination, pro-
duct and quotient (for non-zero denominators) of func-
tions of a Baire class not higher than a is a function in
a Baire class not higher than a. A uniformly convergent
sequence of functions of a Baire class not higher than IX
has as limit a function in a Baire class not higher than
IX. Necessary and sufficient conditions have been esta-
blished for a sequence of functions in a Baire class not
higher than IX to converge to a function in a Baire class
not higher than IX [4]. The union of all dense-in-itself
sets MeA is called the kernel of the topological space A.
If A is a complete space with a non-empty dense-in-
itself kernel, none of the Baire classes is emr)y [2]. The
set of Baire functions coincides with the set of Borel-
measurable functions (cf. Borel measure), and for this
reason they are all Lebesgue-measurable (cf. Lebesgue
measure). A Lebesgue-measurable function g: A ~ R is
equivalent to a Baire function of a Baire class not
higher than two [3]. Baire, who considered functions
defined on Rn (mainly on R), made the most detailed
study of functions of the first class. He showed that a
necessary and sufficient condition for a discontinuous
function to belong to the first class is the existence of a
point of continuity of the induced function on each
perfect set (Baire's theorem). This theorem is applicable
to a function f: A ~R if A has the Baire .property [2].
The concept of a Baire function can be naturally gen-
eralized to functions cp: A ~ Y, where Y is an arbitrary
metric space.
326
References
[1] BAlRE, R.: Lefons sur les fonctions discontinues, professees au
college de France, Gauthier-Villars, 1905.
[2] HAUSDORFF, F.: Set theory, Chelsea, reprint, 1978.
[3] NATANSON, I.P.: Theorie der Funktionen einer reellen
Veriinderlichen, H. Deutsch, Frankfurt a.M., 1961 (translated
from the Russian).
[4) GAGAEFF, B.M.: 'Sur les suites convergentes de fonctions
mesurables B', Fund. Math. 18 (1932), 182-188.
I.A. Vinogradova
Editorial comments. A subset of a topological space is
called dense-in-itse/f if it has no relatively-isolated points. A
modern English reference for the notion of Baire classes is
[A 1]. For an introduction to the topic and some applications,
cf. [A2].
References
[A1] Rooy, A.C.M. VAN and SCHlKHOF, W.H.: A second course
on real functions, Cambridge Univ. Press, 1982.
[A2] BOAS, R.P., JR.: A primer of real functions, Math. Assoc.
Amer., 1981.
AMS 1980 Subject Classification: 26A21 , 54C20,
04A15
BAIRE PROPERTY of a set A in a topological space -
A property analogous to the measurability of a set. A
set A has the Baire property if there exists an open set
G such that the differences A \ G and G \ A are sets of
the first category according to Baire (d. Category of a
set); the term 'open' may be replaced by 'closed'.
There are other equivalent definitions; thus, for
instance, a set has the Baire property if it is the union
of a set of type Go and a set of the first category. The
class of sets with the Baire property is closed under the
operations of complementing, taking countable unions
and taking countable intersections. For an example of
a set which does not have the Baire property, see [1].
References
[I) KURATOWSKl, K.: Topology, Acad. Press, 1966-1968 (translated
from the French). Paragraph 40.
VA. Skvortsov
Editorial comments. A set with the' Baire property is
often called a Baire set or an almost-open set. A set of the
first category is often called meager.
References
[A 1] CECH. E. Topological spaces, Wiley, 1966.
. AMS 1980 Subject Classification: 54C32, 54C50
BAIRE SET in a locally compact Hausdorff space X -
A set belonging to the a-ring generated by the class of
all compact sets in X that are Go-sets. A Baire set
serves to define the concept of a Baire-measurable func-
tion. In all classical particular cases in which measure
theory is developed in topological spaces, e.g. in
Euclidean spaces, the concept of a Baire set coincides
with that of a Borel set.

References
[I) HALMOS, P.: Measure theory, v. Nostrand, 1950.
VA. Skvortsov
AMS 1980 Subject Classification: 28A05, 54C50
BAIRE SPACE - 1) Any space in which the Baire
theorem on complete spaces is valid.
2) The metric space the points of which are infinite
sequences {ni} = { n 1 , n 2, . . . } of natural numbers, and
the distance is given by the formula:
I
p( {n,}, {m,}) = k;'
where ko is the first natural number k for which
nk=l=mk' This is a complete metric separable zero-
dimensional space containing the topological image of
any metric separable zero-dimensional space.
P. S. A leksandrov
AMS 1980 Subject Classification: 54E32
BAIRE THEOREM - 1) Baire's theorem on complete
spaces: Any countable system of open and everywhere-
dense sets in a given complete metric space has a non-
empty, and even an everywhere-dense, intersection in
this space. An equivalent formulation is the following:
A non-empty complete metric space cannot be
represented as a countable sum of its nowhere-dense
subsets. Stated by R. Baire [I].
References
[lJ BAIRE, R.: Ann. Mat. Pura Appl. 3 (1899),67.
P.s. Aleksandrov
Editorial comments. This theorem is also known as the
Baire category theorem (cf. [A1], p. 200).
References
[A 1] KELLEY, J.L.: General topology, van Nostrand, 1955.
2) Baire's theorem on semi-continuous junctions: Let A
be a subset of a metric space M, and let j: M ~R. The
condition: For any number a the set
{x: xEA, j(x)~a} (or, respectively,
{x: xEA, j(x)';:;;a}) is closed in A, is necessary and
sufficient for j (x) to be semi-continuous from above
(Oi, respectively, from below) on A. Demonstrated by
R. Baire for j: R~R [1]. It follows from this theorem
that semi-continuous functions belong to the first Baire
class (cf. Baire classes). A stronger theorem is valid: A
function that is semi-continuous from above (from
below) and that does not assume the value + CXJ ( - (0)
is the limit of a monotone non-increasing (non-
decreasing) sequence of continuous functions.
References
[I] R.: La tMorie desfonctions discontinues, Gauthier-
BAIRE,
Villars, 1905.
[2) NATANSON, LP.: Theorie der Funktionen einer reellen
Veriinderlichen, H. Deutsch, Frankfurt a.M., 1961 (translated
from the Russian). I .A . V'znogra dova
BALANCED SET
Editorial comments. A function in the first Baire class is
also called a Baire function.
References
[A1] RUDIN, W.: Principles of mathematical analysis, McGraw-Hili,
1964.
AMS 1980 Subject Classification: 54E32, 54C50,
26A21
BALANCED MODULE - A module M such that the
natural ring homomorphism cp: R ...... EndEndRM M, where
M is regarded as a right module over EndR M, defined
by cp(r)(m)=mr for any rER and mEM, is surjective.
A module P over a ring R is a generator of the category
of R-modules if and only if P is balanced as an R-
module, projective and finitely generated as an
EndR P-module.
References
[I] FAITH, c.: Algebra: rings, modules and categories, 1-2, Springer,
1973-1976.
L.A. Skornyakov
AMS 1980 Subject Classification: 16A53
BALANCED RING, on the left (right) - A ring over
which all left (right) modules are balanced. A ring is
balanced on the left if and only if all its quotient rings
are QF-l-rings, that is if all the exact left modules over
it are balanced. In particular, a ring is balanced if all
its quotient rings are quasi-Frobenius. Every balanced
ring can be split into a direct sum of a uniserial ring
and rings of matrices over local rings of a special type.
Every balanced ring is semi-perfect. A Noetherian bal-
anced ring is an Artinian ring.
References
[I] Itogi Nauk. i Tekhn. Algebra Topol. Geom. 19 (1981),31-134.
[2) FAITH, c.: Algebra, 1-2, Springer, 1973-1976.
L.A. Skornyakov
AMS 1980 Subject Classification: 16A36
BALANCED SET - A set U in a real or complex vec-
tor space X such that x E U and 1,\ I .;:;; I imply Ax E U.
An example of a balanced set is given by the unit ball
in a normed vector space and, generally, by a neigh-
bourhood U of zero in a base of neighbourhoods of
zero in a topological vector space. These neighbour-
hoods of zero are moreover absorbing, i.e. such that for
any x EX there exists an 0:>0 such that x E'\U for
I ,\ I ~a. If U is a convex, absorbing and balanced set,
then the functional Pu(x)=inf{ 1,\ I: x E'\U} is a
semi-norm, i.e. it has the properties
Pu(x +y) ,.;; Pu(x)+pu(y), PU(Ax) = 1,\ Ipu(x).
A balanced set is also called centred.
References
[I] KANrOROVICH, L.V. and AKILOv, G.P.: Functional analySis,
Moscow, 1977 (in Russian).
VI. Sobolev
327
BALANCED SET
Editorial comments. The functional Pu mentioned above
is also called the Minkowski functional of the convex,
absorbing and balanced set U.
References
[A1] YOSIDA, K.: Functional analysis, Springer, 1980.
AMS 1980 Subject Classification: 46AXX
BALAYAGE METHOD - A method for solving the
Dirichlet problem for the Laplace equation, developed
by H. Poincare ([1], [2], see also [4]), which will now be
described. Let D be a bounded domain of the
Euclidean space Rn , n;:' 2, let f = aD
be the boundary
of D. Let 8y be the Dirac measure concentrated at the
point y ED, let U(x; 8y ) be the Newton potential of the
measure 8y for n;:' 3, or the logarithmic potential of the
measure 8y if n = 2. A balayage (or sweeping) of the
measure 8y from the domain D to the boundary f is a
measure /3y on f whose potential U(x; /3y ) coincides
outside D with U(x; 8y ) and is not larger than U(x; 8y )
inside D; this measure /3y is unique and coincides with
the harmonic measure on f for the point y ED. The
balayage of an arbitrary positive measure, concentrated
on D, is defined in a similar manner. If D is a sphere,
the density of the mass distribution /3Y' i.e. the deriva-
tive of the measure /3y, is identical with the Poisson
kernel (cf. Poisson integral). In general, if the boundary
r is sufficiently smooth, the measure /3y is absolutely
continuous, and the density of the mass distribution /3y
coincides with the normal derivative of the Green func-
tion for D. The measure /3y serves to write down the
solution of the Dirichlet problem as the so-called for-
mula of de la Val!ee-Poussin:
U(y) = /f (x) d/3y(x),
where f (x) is a function defined on f.
In his original publication on the balayage method,
Poincare began by demonstrating the geom) trical con-
struction of the process for a sphere. Then, Lasing him-
self on Harnack's theorems (cf. Harnack theorem) and
on the fact that it is possible to exhaust the domain D
r
by a sequence of spheres {Bd = I, he constructed an
infinite sequence of potentials {Un} ~= I in which each
potential Un + I is obtained from the preceding one, Un,
by the balayage method of moving the masses from the
domain U: = 1Bk to its boundary, and which reduces
to solving the Dirichlet problem for a sufficiently
smooth domain D (for a detailed discussion of the con-
ditions of applicability of the balayage method, see [3]).
In modern potential theory [5], [6] the balayage prob-
lem is treated as an independem problem, resembling
the Dirichlet problem, and it turns out that the balay-
aged measure can be considered on sets of a general
nature. For instance, the balayage problem in its sim-
328
plest form is to find, fora given mass distribution J.L
inside a closed domain 15, a mass distribution p on
f= aD such that the potentials of both distributions
coincide outside D. If the boundary f is smooth, the
solution of the balayage problem for J.L will be an abso-
lutely continuous measure P. Its density, or the
derivative p' 0'), y Ef, may be written down in terms of
the Green function G(x, y) of the domain D in the
form
p'(y) = j t,
aGe Y) dp.(x),
ny
p.;;;'O, yEf, (*)
where aG / any is the derivative of G (x, y) in the direc-
tion of the interior normal to f at the point y Ef.
Inside the domain D the potentials satisfy the inequal-
ity U(x; p)";;; U(x; p.), i.e. balayage inside the domain
results in a decrease of the potential. If p.=8x is the
Dirac measure at the point x ED, formula (*) yields
p'(y)=aG(x,y)/any , i.e. the normal derivative of the
Green function is the density of the measure obtained
by balayage of the unit mass concentrated at the point
x ED. Generalization of formula (*) yields an expres-
sion for the balayaged measure p(E) of an arbitrary
Borel set E cf for an arbitrary domain D:
pCE) = jw(x;E, D)dp.(x),
where w(x; E, D) is the harmonic measure of E with
respect to the domain D at the point x.
If K is an arbitrary compact set in Rn and p. is a
bounded positive Borel measure, the balayage (or
sweeping) of the measure p. onto the compact set K is a
measure p on K such that U(x; p)";;;; U(x; p.) every-
where, and such that quasi-everywhere on K, i.e. with
the possible exception of a set of points of exterior
capacity zero, U(x; p)= U(x; p.). Such a formulation of
the balayage problem, which is more general than
balayage from a domain, may also be extended to
potentials of other types, e.g. Bessel potentials or Riesz
potentials (cf. Bessel potential; Riesz potential). Balay-
age of measures onto arbitrary Borel sets K is also con-
sidered.
The problem of balayage for superharmonic func-
tions (cf. Superharmonic function) has been similarly
formulated. Let v be a non-negative super harmonic
function on a domain D eRn. The balayage of the
function v onto a compact set K C 15 is the largest
superharmonic function B~ such that 1) its associated
measure is concentrated on K; 2) B~";;;;v everywhere;
and 3) B~ = v quasi-everywhere on K.
In abstract potential theory (cf. Potential theory,
abstract) the balayage problem in both its formulations
is solved for sets K in an arbitrary harmonic space X,
i.e. in a locally compact topological space X which per-
mits the isolation of an axiomatically defined sheaf of
harmonic functions. This axiomatic approach makes it
possible to consider the balayage problem for potentials

connected with partial differential equations of a more
general nature [7]. For the balayage method in stochas-
tics cf. [8].
References
[1) POINCARE, H.: 'Sur Ies equations aux derivees partielles de Ia
physique mathematique', Amer. J. Math. 12, no. 3 (1890), 211-
294.
(2) POINCARE, H.: Theorie du potentiel Newtonien, Paris, 1899.
(3) VALLEE-POUSSIN, Cu.J. DE LA: Le potentiellogarithmique,
balayage et representation confonne, Gauthier-Villars, 1949.
(4) SRETENSKIi, L.N.: Theory of the Newton potential, Moscow-
Leningrad, 1946 (in Russian).
(5) LANDKOF, N.S.: Foundations of modem potential theory,
Springer, 1972 (translated from the Russian).
(6) BRELOT, M.: Elements de la tMorie c/assique du potentiel, Sor-
bonne Univ., Paris, 1965.
(7) CONSTANTINESCU, C. and CORNEA, A.: Potential theory on
harmonic spaces, Springer, 1972.
(8) MEYER, P.A.: Probability and potentials, Blaisdell, 1966.
E.D. Solomentsev
Editorial comments. 8alayage is also referred to as
sweeping of a measure. A classic reference for problems in
potential theory related to Green functions is [A 1].
In probabilistic potential theory the swept measure P.s on
f=aD of a probability measure p. concentrated on D turns
out to be the distribution of a standard Brownian motion on
D, which has initial distribution p., at the moment of first hit-
ting f.
Another link with probabilistic potential theory is provided
by the fact that, for each sufficiently nice harmonic space,
there exists a Hunt process (Xt ) whose excessive functions
are the positive hyper-harmonic functions. If PK denotes the
hitting distribution of a compact set K, then PKu=B~ for
positive superharmonic functions u, and the balayage of a
measure p. on K is given by P.PK' Therefore, the notion of
balayage of a function or a measure can also be defined in
terms of the potential kernel of a semi-group of kernels, see
[A3].
References
[A1) TSUJI, M.: Potential theory in modern function theory, Chel-
sea, reprint, 1975.
[A2] DooB, J.L.: Classical potential theory and its probabilistic
counterpart, Springer, 1984.
[A3] DELLACHERIE, C. and MEYER, P.A.: Probabilites et potentiel,
1-2, Hermann, 1975-1983.
AMS 1980 Subject Classification: 31825, 31815,
31820,60J45
BALL - A set vn of points x in a Euclidean space
En at a distance from a given point Xo (the centre of
the ball) less than (an open ball von), or not greater
than (a closed ball V") a quantity R (the radius of the
ball), i.e.
von (V') = {XEE": p(x, xo)<R (";;;;;R)}.
A ball Vi is a line segment, V2 is a disc, vn for n > 3 is
sometimes called a hyperball. The boundary (surface) of
a ball is a sphere.
The volume of a ball is
BANACH ALGEBRA
'11'/2 s"-I
V" = f(n /2+ 1) R" = -n- R ,
where sn -I is the surface of the boundary sphere and
f is the gamma-function: f(n + l)=n!,
f(n + I / 2)=2- n (l'3 ... (2n -1)Jn.l /2. In particular,
VI = 2R
,
V 2 = wR2
,
V3 = .i
3 wR 3' V4 = ~2R4 .
With the increase of the dimension, the volume of a
ball 'concentrates' at its surface:
vn(R2)- V"(R I )
V"(R 2) ~ 1, RI <R 2, n~oo.
A ball is the simplest geometrical figure. Its topology
is trivial. Among all bodies of an equal volume, a ball
has minimal surface, and among all bodies of an equal
surface, it has maximal volume.
In exactly the same manner a ball can be defined in
a metric space; however, in this case it need not be, for
example, strictly convex, its surface may have non-
smooth points, etc., that is, it may have all phenomena
characteristic of arbitrary convex bodies.
Unlike a finite-dimensional ball, an inJinite-
dimensional ball, being the direct limit of a sequence of
balls of succesive dimensions imbedded in one another,
does not have a compact closure. On the contrary, the
compactness of a ball in a topological vector space
indicates the finite dimensionality of the latter.
For references, see Sphere.
I.S. Sharadze
AMS 1980 Subject Classification: 51 N20, 46805
BANACH ALGEBRA - A topological algebra A over
the field of complex numbers whose topology is defined
by a norm which converts A into a Banach space, the
multiplication of the elements being separately continu-
ous for both factors. A Banach algebra is said to be
commutative if xy = yx for all x, y EA (cf. Commutative
Banach algebra). A Banach algebra is said to be an
algebra with a unit if A contains an element e such that
ex =xe =x for any x EA. If a Banach algebra has no
unit, a unit may be adjoined, i.e. it is possible to con-
s_truct a Banach algebra A with a unit element such that
A contains the initial algebra A as a closed sub algebra
of codimension one. In any Banach algebra A with a
unit element e it is possible to change the norm for an
equivalent one so that in the new norm the relation-
ships II ab I , ; ; ; I a II I b II, I e I = I are valid. In what
follows it is, as a rule, assumed that the algebra does
contain a unit and that it satisfies the norm conditions
given above.
Examples. I) Let X be a compact topological space
and let C (X) be the set of all continuous complex-
valued functions on X. C (X) will then be a Banach
algebra with respect to the usual operations, with norm
329
BANACH ALGEBRA
11I11 = m:x III
2) The set of all bounded linear operators on a
Banach space forms a Banach algebra with respect to
the usual operations of addition and multiplication of
linear operators with the operator norm.
3) Let V be a bounded domain in n-dimensional
complex space en. The set of bounded holomorphic
functions on V is a Banach algebra with respect to the
usual operations, with the natural sup-norm:
I I II = sttP II I
This Banach algebra contains the closed subalgebra of
bounded holomorphic functions on V that have a con-
tinuous extension to the closure of V. The simplest
example is the algebra of functions that are continuous
in the disc I z I ,;;;; 1 and analytic in the disc I z I < 1.
4) Let G be a locally compact group and let L] (G)
be the space (of equivalence classes) of all functions
that are measurable with respect to the Haar measure
on G and that are absolutely integrable with respect to
this measure, with norm
I I I = I Il(g) I dg
G The set of elements invertible from the right (from the
(left Haar integral).
If the convolution operation
(II *h)(h) = III (g)h(g-I h)dg
G morphic to the field of complex numbers (the
is considered as the multiplication in L] (G), then
L] (G) becomes a Banach algebra; if G is an Abelian
locally compact group, then the Banach algebra L] (G)
is commutative. The Banach algebra L](G) is said to
be the group algebra of G. The group algebra L] (G)
has a unit (with respect to the convolution) if and only
if G is discrete.
If G is commutative it is possible to construct a faith-
ful representation of L] (G), given by the Fourier
transform of each function f EL] (G), i.e. bJ the func-
tion
I(x) = Ix(g)/(g)dg,
G
on the character group G of p. The set of functions f(~)
A A
form;; a certain algebra A (G) of continuous functions
on G (with respect to the ordinary pointwise opera-
tions), called the Fourier algebra of the locally compact
Abelian group G. !n particular, if G is the group of
integers Z, then A (Z) is the algebra of continuous func-
tions on the circle which are expandable into an abso-
lutely convergent trigonometric series.
5) Let G be a topological group. A continuous
complex-valued function f (g) on G is said to be almost
periodic if the set of its shifts f(gog), go E G, forms a
compact family with respect to uniform convergence on
G. The set of almost-periodic functions forms a com-
330
mutative Banach algebra with respect to the pointwise
operations, with norm
II I II = sup
get;
1/(g) I
6) The skew-field of quaternions does not form a
Banach algebra over the field of complex numbers,
since the product of elements of a Banach algebra A
should be compatible with multiplication by numbers:
For all AE e and x, y EA the equation
>..(xy) = (Ax)y = x(>"y),
must be valid; it is not valid in the field of quaternions
if A=i, x=j,y =k.
Any Banach algebra with a unit is a topological alge-
bra with continuous inverses. Moreover, if (A) is the
set of elements of a Banach algebra A which have a
(two-sided) inverse with respect to multiplication, then
(A) is a topological group in the topology induced by
the imbedding (A)CA. If II e-a I <1, then aE(A),
and
where b = e - a, and the series is absolutely convergent.
left) in A also forms an open set in A.
If in a Banach algebra A all elements have an inverse
(or even a left inverse), then A is isometrically iso
Gel'fand- Mazur theorem).
Since a certain neighbourhood of the unit in a
Banach algebra A consists of invertible elements, the
closure of any non-trivial ideal is again an ideal which
does not coincide with A. In particular, a maximal (left,
right, two-sided) ideal is closed.
An important task in the theory of Banach algebras
is the description of closed ideals in Banach algebras.
The problem can be simply solved in a number of
cases. In the algebra C(X) (cf. Example 1) each closed
ideal has the form 1= {fE C (X): fl y ~O}, where Y is a
closed set in X. If A is the algebra of all bounded linear
operators on a separable infinite-dimensional Hilbert
space, then the ideal of completely-continuous opera-
tors is the only closed two-sided ideal in A.
An element a EA has a left (right) inverse if and only
if it is not contained in any maximal left (right) ideal.
The intersection of all maximal left ideals in A coin-
cides with the intersection of all maximal right ideals;
this intersection is called the radical of the algebra A
and is denoted by RadA. An element ao EA belongs to
RadA if and only if e + aao E(A) for any a EA. Alge-
bras for which RadA =0 are said to be semi-simple.
The algebras C( X) and the group algebras L] (G) are
semi-simple. All irreducible (i.e. not having a non-
trivial invariant subspace) closed sub algebras of the

algebra of all bounded linear operators on a Banach
space are semi-simple.
The resolvent of an element a EA is the function
l>.~a" = (a-l>.e)-I,
defined on the set of all AEC for which a (two-sided)
inverse to a - Ae exists. The domain of existence of the
resolvent contains all points A with I A I > II a II. The
maximal domain of existence of the resolvent is an
open set; the resolvent is continuous on this set and is
even analytic, moreover da}, / dA=aX. In addition,
Hilbert's identity
a"2 -a", = (l>.2-l>.da",a"2
is valid. The complement of the domain of existence of
the resolvent is called the spectrum of the element a and
is denoted by a(a). For each a EA the set a(a) is non-
empty, closed and bounded.
If a, bEA, then the sets a(ab) and a(ba) need not
coincide, but
a(ab) U {O} = a(ba) U {O}.
The number
lal = max Il>.l
"Eo(a)
is called the spectral radius of the element a EA;
Gelfand's formula
I a 1= lim I an III/n,
where the limit on the right-hand side always exists, is
valid. If a ERadA, then I a I =0; the converse is true,
generally speaking, only in commutative Banach alge-
bras whose radical coincides with the set of generalized
nilpotents, i.e. elements a for which I a I = O. In any
Banach algebra the relationships I a k I = I a Ik,
I Aa I = I A I I a I and I a I ~ II a II are true. If A is
commutative, then I ab I ~ I a I I b I and
I a + b I ~ I a I + I b I are valid.
Examples of non-commutative algebras in which gen-
eralized non-zero nilpotents are absent are known.
However, if II a 2 II = I a 112 for any a EA, then the
Banach algebra A is commutative. The condition
II ab II = I ba II for all a, b EA is also sufficient for an
algebra A with a unit to be commutative.
An algebra A is said to be an algebra with involution
if an operation a ----'?a * is defined on A that satisfies the
conditions
(l>.a +ph)' = Xa' +jih', (a')' = a, (ab)' =b' Q',
for all a, bEA, A, IlEe. The mapping a----'?a* is said to
be an involution in A. A linear functional1/; on an alge-
bra A with an involution is said to be positive if
1/;(aa*r;:.O for any aEA. If the linear functional1/; is
positive, then
for all a EA. If the involution in A is an isometry, i.e. if
BANACH ALGEBRA
II a* II = I a II for all aEA, then
t/J(a' a) ~ I/{e) I a' a I
A Banach algebra A with involution is said to be
completely symmetric if e + a * a E (A) for any a EA; A is
said to be a C* -algebra (a completely-regular algebra) if
II a *a II = I a 112 for any a EA. Any C* -algebra is com-
pletely symmetric. Examples of completely-symmetric
algebras include the group algebras L 1(G) of commuta-
tive or compact groups. Examples of C* -algebras
include the algebras C(X) (the involution in C(X) is
defined as transition to the complex conjugate func-
tion) and closed subalgebras of the algebra of bounded
linear operators in a Hilbert space containing both the
operator and the adjoint operator (involution is defined
as transition to the adjoint operator). Any C* -algebra
is isometrically isomorphic (involution being preserved)
with one of these algebras (the Gel'fand- Narmark
theorem). In particular, any commutative C -algebra A
is isometrically isomorphic (involution being preserved)
with one of the algebras C(X) (this theorem includes
the Stone- Weierstrass theorem).
An element a of a Banach algebra with involution is
said to be Hermitian if a* =a. For a Banach algebra
with an involution to be a C* -algebra it is necessary
and sufficient that the condition II e ia II = 1 be fulfilled
for all Hermitian elements a. If, in a Banach algebra
with an involution, sup II e ia II < 00 (upper bound over
all Hermitian elements), then this algebra is topologi-
cally *-isomorphic with a C* -algebra. If, in an arbi-
trary Banach algebra, II e ita II = 1 for all real t for a cer-
tain fixed element a, then II a II coincides with the
spectral radius, i.e. I a I = I a I
The theory of Banach algebras, and of commutative
Banach algebras in particular, has numerous applica-
tions in various branches of functional analysis and in
a number of other mathematical disciplines.
References
[I] BOURBAKI, N.: Elements of mathematics. Spectral theories,
Addison-Wesley, 1977 (translated from the French).
[2] GAMELIN, T.W.: Uniform algebras, Prentice-Hall, 1969.
[3] GUNNING, R.C. and ROSSI, H.: Analytic functions of several
complex variables, Prentice-Hall, 1965.
[4] GEL'FAND, I.M.: 'Normierte Ringe', Mat. Sb. 9 (51), no. I
(1941), 3-24.
[5] GLEASON, A.M.: 'Function algebras', in Proc. Sem. on analytiC
functions, Vol. 2, 1958, pp. 213-226.
[6] HOFFMAN, K.: Banach spaces of analytic functions, Prentice
hall, 1962.
[7] GORIN, E.A.: 'Maximal sub algebras of commutative Banach
algebras with involution', Math. Notes 1, no. 2 (1967), 173-178.
(Mat. Zametki 1, no. 2 (1967),173-178)
[8] DUNFORD, N. and SCHWARTZ, J.T.: Linear operators, 1-3,
Interscience, 1958-1971.
[9] ZELAZKO, W.: Banach algebras, Elsevier, 1973 (translated from
the Polish).
[10] KAPLANSKY, I.: 'Functional analysis', in Surveys in applied
mathematics, Vol. 4. Some aspects of analysis and probability,
Wiley, 1958.
331
BANACH ALGEBRA
[11] LoOMIS, L.H.: An introduction to abstract harmonic analysis, v.
Nostrand, 1953.
[12] NAiMARK, M.A.: Normed rings, Reidel, 1984 (translated from
the Russian).
[13] Some questions in approximation theory, Moscow, 1963 (in Rus-
sian; translated from the English).
[14] RICKART, C.E.: General theory of Banach algebras, Nostrand,
1960.
[15] ROYDEN, H.L.: 'Function algebras', Bull. Amer. Math. Soc. 69,
no. 3 (1963), 281-298.
E.A. Gorin
Editorial comments. Gel'fand's formula is also called the
spectral radius formula.
References
[A1] KADlSON, R.V. and RINGROSE, l.R.: Fundamentals of the
theory of operator algebras, 1, Acad. Press, 1983.
AMS 1980 Subject Classification: 46H05, 46J05,
43A10
BANACH ANALYTIC SPACE An infinite-
dimensional generalization of the concept of an analytic
space, which arose in the context of the study of defor-
mations of analytic structures (d. Defonnation). Here,
the local model is a Banach analytic set, i.e. a subset
p.( U, F, j) = f- 1(0) of an open set U in a Banach space
E over C, where f: U ~F is an analytic mapping into
the Banach space F. As distinct from the finite-
dimensional case, not one structure sheaf, but a set of
sheaves <1>( W), where W is an open set in an arbitrary
Banach space G, is defined on the local model. <1>( G) is
defined as the quotient of the sheaf of germs of analytic
mappings U ~G by the sub sheaf of germs of mappings
of the type x~cp(x)f(x), where </>: U~Hom(F, G) is a
local analytic mapping, while <1>( W) c <1>( G) is generated
by mappings which assume values in W. The sheaves
<1>( W) define a functor from the category K of open sets
in Banach spaces and their analytic mappings into the
category of sheaves of sets on f- 1(0).
A topological space X with a functor from the
category K into the category of sheaves of s) ts in X in
which all points have neighbourhoods isomorphic to
some local model, is said to be a Banach analytic space.
Complex-analytic spaces form a complete sub-
category in the category of Banach analytic spaces. A
Banach analytic space is finite-dimensional if each one
of its points x has a neighbourhood that is isomorphic
to a model p.( U, F, j) and for which there exists a map-
ping g: U ~ U inducing an automorphism of the model
and having a completely-continuous differential dgx [1].
A second special case of a Banach analytic space is a
Banach analytic manifold, i.e. an analytic space that is
locally isomorphic to open subsets of Banach spaces.
An important example is the manifold or
linear sub-
spaces of a Banach space over C that are closed and
admit closed complements.
Finitely-defined Banach analytic sets, i.e. models of
332
the type p.( U, Ck , f), have local properties which
correspond to classical properties: primary decomposi-
tion, Hilbert's Nullstellen theorem, the local description
theorem, etc., are all applicable [2].
References
[I] DouADY, A.: 'Les problemes des modules pour les sous-espace
analytique compacts d'une espace analytique donne', Ann. Inst.
Fourier (Grenoble) 16, no. 1 (1966), 1-95.
[2] Sous-ensembles analytiques d'une variete banachique complexe,
Springer, 1970.
D.A. Ponomarev
AMS 1980 Subject Classification: 46BXX, 32CXX
BANACH INDICATRIX, multiplicity function, of a
continuous function y = f (x), a";;;;x,,;;;;b - An integer-
valued function N(y,j), -oo<y<+oo, equal to the
number of roots of the equation f (x) = y. If, for a given
value of y, this equation has an infinite number of
roots, then
N(Y,f) = +00,
and if it has no roots, then
N(y,f) = O.
The function N(y,f) was defined by S. Banach [1] (see
also [2]). He proved that the indicatrix N (y, f) of any
continuous function f(x) in the interval [a, b] IS a
function of Baire class not higher than 2, and
+00
f N(y, f)dy,
b
V(f) = (*)
a -00
where V;;(f) is the variation of f (x) on [a, b]. Thus,
equation (*) can be considered as the definition of the
variation of a continuous function f (x). The Banach
indicatrix is also defined (preserving equation (* for
functions with discontinuities of the first kind [3]. The
concept of a Banach indicatrix was employed to define
the variation of functions in several variables [4], [5].
References
[1] BANACH, S.: 'Sur les lignes rectifiables et 1es surfaces dont
l'aire est finie', Fund. Math. 7 (1925),225-236.
[2] NATANSON, J.P.: Theorie der Funktionen einer reellen
Veriinderlichen, H. Deutsch, Frankfurt a.M., 1961 (translated
from the Russian).
[3] LOZINSKI), S.M.: 'On the Banach indicatrix', Vestnik Len-
ingrad. Univ. Math. Mekh. Astr. 7, no. 2, 70-87 (in Russian).
[4] KRONROD, A.S.: 'On functions of two variables', Uspekhi Mat.
Nauk 5, no. 1 (1950),24-134 (in Russian).
[5] VITUSHKIN, A.G.: On multi-dimensional variations, Moscow,
1955 (in Russian).
B.I. Golubov
Editorial comments. More generally, for any mapping
t: X---'> Y define NCy, t) analogously. Then, let X be a separ-
able metric space and let tCA) be /L-measurable for all Borel
subsets A of X. Let ~Cs)=/L(fCs) for Sex and let>/; be the
measure on X defined by the Caratheodory construction
from ~. Then
>/;CA) = f
N(f)A) d/Ly
for every Borel set Aex. Cf. [A1], p. 176 ft. For significant
extension of ("), cf. [A2].

Editorial comments.
References
[A 1] FEDERER, H.: Geometric measure theory, Springer, 1969.
[A2] FEDERER, H.: 'An analytic characterization of distributions
whose partial derivatives are representable by measures',
Bull. Amer. Math. Soc. 60 (1954),339.
AMS 1980 Subject Classification: 26A15
BANACH LATIICE, - A vector lattice that is at the
same time a Banach space with a norm which satisfies
the monotonicity condition:
I x I ~ Iy I ~ II x I ~ I y II
A Banach lattice is also called a KB-lineal, whereas an
arbitrary normed lattice, i.e. a vector lattice with a
monotone norm, is called a KN-lineal. When complet-
ing a normed lattice in norm, the order relations may
be extended to the resulting Banach space so that it
becomes a Banach lattice. If it is possible to introduce
in a lattice a Banach topology which converts it to a
Banach lattice, such a topology is unique. The simplest
example of a Banach lattice is the space C(Q) of con-
tinuous functions on an arbitrary compact topological
space Q with the natural (pointwise) order and with the
ordinary (uniform) norm. Other examples of Banach
lattices include Lp spaces and Orlicz spaces (cf. Orlicz
space). In Banach lattices convergence in norm is (*)-
convergence for convergence with a regulator. This is
not true of normed lattices.
An important special case is a Banach lattice of
bounded elements. If a lattice X contains a strong unit
1, i.e. if for each x EX there exists a ,\ such that
I x I :S;;;,\I, then the smallest ,\ for which this inequality
is valid is taken as II x II. The normed lattice thus
obtained is called a normed lattice of bounded elements;
if it is complete in norm, it is called a Banach lattice of
bounded elements. In a Banach lattice (and even in a
normed lattice) of bounded elements convergence in
norm is identical with convergence with a regulator,
while the boundedness of a set of elements in norm is
identical with order boundedness. If a normed lattice of
bounded elements is conditionally a-complete, it is
complete in norm.
The space C(Q) is a Banach lattice of bounded ele-
ments in which the function x(q) 1 is taken as the
unit. For any Banach lattice X of bounded elements
there exists a compact Hausdorff space Q such that X
is algebraically and lattice isomorphic to the space
C (Q). This is an abstract characterization of the
Banach lattice of continuous functions on a compact
Hausdorff space.
In any normed lattice an additive functional that is
continuous in norm is regular and, moreover, is
representable as the difference of two additive function-
als which are continuous in norm. In a Banach lattice
BANACH - MAZUR FUNCTIONAL
each positive additive functional is continuous in norm,
which means that the classes of regular functionals and
additive functionals which are continuous in norm
coincide. The space X' which is dual in the sense of
Banach to the normed lattice X is a conditionally com-
plete Banach lattice. In a normed lattice the
Hahn - Banach theorem may be strengthened as fol-
lows: For any Xo >0 there exists a positive additive
functional f, which is continuous in norm, such that
f(xo) = II Xo II, I f II = 1.
References
(1) VULIKH, B.Z.: Introduction to the theory of partially ordered
spaces, Wolters-Noordhoff, 1967 (translated from the Russian).
[2) DAY, M.M.: Normed linear spaces, Springer, 1958.
B.Z. Vulikh
Editorial comments. Cf. Riesz space for the notion of
'convergence with a regulator'. The terminology in the arti-
cle above is from [1], and is not in general use in the
Western literature.
A vector lattice (i.e. a real linear (partially) ordered vector
space that is at the same time a lattice) is usually called a
Riesz space. A norm II . lion a Riesz space L is called a
Riesz norm if
I t I ~ I g I ~ I t II ~ II g II
whenever t, gEL (called the monotonicity condition,
above). Here
ItI = (- t)V(f).
A normed Riesz space is understood to mean a Riesz space
provided with a Riesz norm. The definition of a Banach lat-
tice can now be rephrased as follows: A Banach lattice is a
normed Riesz space that is complete in the norm.
All (Riesz) norms on a Banach lattice are equivalent.
For a Banach lattice L the order dual space (cf. Riesz
space) and norm dual space coincide (this is not true for
normed Riesz spaces in general). A Banach lattice, required
to be complete in norm, need not be Oedekind complete.
However, any Banach lattice L (even, every normed Riesz
space) can be regarded as a (Riesz) subspace of a
Oedekind-complete Banach lattice (the bi-dual (in norm) of
L, cf. Banach space).
References
[A1] SCHAEFER, H.H.: Banach lattices and positive operators,
Springer, 1974.
[A2] LUXEMBURG, W.A.l. and ZAANEN, A.c.: Riesz spaces, 1,
North-Holland, 1971.
[A3] LINDENSTRAUSS, 1. and TZAFRlRl, L.: Classical Banach
spaces, 2, Springer, 1979.
[A4] ZAANEN, A.c.: Riesz spaces, 2, North-Holland, 1983.
AMS 1980 Subject Classification: 46A40, 46830
BANACH - MAzUR FUNCfIONAL,
Banach - Mazur operator - A concept of a computable
functional (operator), proposed by S. Banach and S.
Mazur [l] and concerning the computability of a func-
tional (operator) from a set M 1 into a set M 2 as its
property to transform any computable sequence of e1e-
333
BANACH - MAZUR FUNCTIONAL
ments of M 1 into a computable sequence of elements
of M2 (d. Computable function).
Let R be the set of all one-place general recursive
functions (d. General recursive function). A functional
<I> defined on R and assuming natural numbers as
values is called computable according to
Banach - Mazur, or a Banach - Mazur functional, if for
any two-place general recursive function g there exists a
general recursive function f such that
j(n) = <I>(g(n, m)).
(Here g is considered to be a function of m for any
constant n.) All general recursive functionals and
everywhere-defined effective functionals (d. Construc-
tive metric space) are Banach - Mazur functionals. On
the other hand, an example of a Banach - Mazur func-
tional which is not identical with any general recursive,
and hence neither with any effective, functional is
known [2]. An important property of Banach - Mazur
functionals is their continuity [1]: The values of such a
functional on any general recursive function is defined
only by a finite number of values of this function.
The concept of computability just outlined can be
extended to functions of a real variable. Let C be the
set of computable sequences of computable real
numbers; each sequence {ak} E C is defined by two gen-
eral recursive functions f and g such that for all nand
k la- (k, n)-g(k,
k n+l
n) I < _1_.
n+l
A function cf> of a real variable is said to be computable
according to Banach - Mazur (the set of such functions
is denoted by 91) if for any sequence {ak} from C the
sequence {cf>( ak)} also belongs to C. All functions cf> E 91
are continuous at all computable points [1]. For
instance, sgn x f/:. 91. The question as to whether all
functions from 91 are computably continuous is still
(1977) open. The set 91 is closed with resr~ct to the
sequence of operations used in analysis, so) that com-
putable analysis can be successfully developed on this
basis [1].
References
[1] MAZUR, S.: Computahle analysis. PWN. 1963.
[2] FRIEDBERG. R.M.: '4-quantifier completeness: A Banach-
Mazur functional not uniformly partial recursive', Bull. Acad.
Polan. Sci. Ser. Sci. Malh., ASlr. Phys. 6, no. I (1958 J. 1-5.
[3] MARKOV, A.A.: 'On constructive functions'. Trud)' Mal. Jnsl.
Sleklov. 52 (1958), 315-348 (in Russian).
[4] ROGERS, JR .. H.: Theory of recursive junclions and ejfeclive
compulahilill'. McGraw-Hill, 1967.
B.A. Kus ncr
Editorial comments. For a more recent exposition of
recursive analysis, see [A1].
References
[A1] AIIERTH. O. Computable analysis. McGraw-Hili. 1980.
AMS 1980 Subject Classification: 03020, 03F60
334
BANACH MODULE (left) over a Banach algebra A -
A Banach space X together with a continuous bilinear
operator m:A XX~X defining on X the structure of a
left module over A in the algebraic sense. A right
Banach module and a Banach bimodule over A are
defined in an analogous manner. A continuous
homomorphism of two Banach modules is called a
morphism. Examples of Banach modules over A
include a closed ideal in A and a Banach algebra B :JA.
A Banach module over A that can be represented as a
direct factor of Banach modules A + E, (where A + is
A with an added unit and E is a Banach space and
mea, bx)=abx) is called projective. Cf. Topologi-
cal tensor product.
References
[1] R.iEFFEL, M.A.: 'Induced Banach representations of Banach
algebras and locally compact groups', J. Funct. Anal. 1 (1967),
443-491. A .1.v a. lGh eIemsk'II v
AMS 1980 Subject Classification: 46825
BANACH SPACE, B-space - A complete normed vec-
tor space. The function spaces introduced by D. Hil-
bert, M. Frechet and F. Riesz between 1904 and 1918
served as the starting point for the theory of Banach
spaces. It is in these spaces that the fundamental con-
cepts of strong and weak convergence, compactness,
linear functional, linear operator, etc., were originally
studied. Banach spaces were named after S. Banach
who in 1922 began a systematic study of these spaces,
based on axioms introduced by himself, and who
obtained highly advanced results.
The theory of Banach spaces developed in parallel
with the general theory of linear topological spaces (cf.
Linear topological space). These theories mutually
enriched one another with new ideas and facts. Thus,
the idea of semi-norms, taken from the theory of
normed spaces, became an indispensable tool in con-
structing the theory of locally convex linear topological
spaces. The ideas of weak convergence of elements and
linear functionals in Banach spaces ultimately evolved
to the concept of weak topology. The theory of Banach
spaces is a thoroughly studied branch of functional
analysis, with numerous applications in various
branches of mathematics - directly or by way of the
theory of operators.
The problems involved in Banach spaces are of dif-
ferent types: the geometry of the unit ball, the
h geometry of subspaces, the linear topological classifica-
tion, series and sequences in Banach spaces, best
approximations in Banach spaces, functions with values
in a Banach space, etc. Regarding the theory of opera-
tors in Banach spaces it should be pointed out that
many theorems are directly related to the geometry and
the topology of Banach spaces.
 BANACH SPACE

Examples. The Banach spaces encountered in Lp(S;~, JL). It is the space of Lebesgue-measurable
analysis are mostly sets of functions or sequences of functions, summable of degree p, with the norm
numbers which are subject to certain conditions.
I) ~,p;;;;'1, is the space of numerical sequences II x I = [[I xes) IPdsfIP.
x={~n} for which
00
13) AP is the Bohr space of almost-periodic func-
~ I ~n IP < 00, tions, with the norm
n=1

with the norm Ilxll= sup Ix(t)l.

-00<1<::+00

The spaces qa, b], C[a, b], Lp[a, b], c, ~ are

separable; the spaces M[ a, b], m, AP are non-
separable; C[K] is separable if and only if K is a com-
2) m is the space of bounded numerical sequences
pact metric space.
with the norm
I x I = sup
n
I ~n I A (closed linear) subspace Y of a Banach space, con-
sidered apart from the enveloping space X, is a Banach
3) c is the space of convergent numerical sequences space. The quotient space X / Y of a normed space by
with the norm a subspace Y is a normed space if the norm is defined
I x II = sup I ~n I
n
as follows. Let Y 1 = X 1 + Y be a coset. Then
4) Co is the space of numerical sequences which
converge to zero with the norm
II YII = YEY
infII X1+Y II
If X is a Banach space, then X / Y is a Banach space
II x II = max I ~n I
n
as well. The set of all continuous linear functionals
5) qa, b] is the space of continuous functions defined on the normed space X, with the norm
x=x(t) on [a, b] with the norm
11/11 = sUP I ((x) I , x~O,
IxI = max
a:<;;;,ts;,b
I x(t) I XEX I x II
is said to be the dual space of X, and is denoted by X'.
6) C[K] is the space of continuous functions on a
It is a Banach space.
compactum K with the norm
Banach spaces satisfy the Hahn - Banach theorem on
II x II = max I x(t) I
IEK
the extension of linear functionals: If a linear func-
tional is defined on a subspace Y of a normed space X,
7) C[a, b] is the space of functions with continuous
it can be extended, while preserving its linearity and
derivatives up to and including the order n, with the
continuity, onto the whole space X. Moreover, the
norm
I x II = k~=(j1";;I,,;;b
max I X(k)(t) I
extension can be made to have the same norm:
IIfllx = sUP I (x) I = IIflly = sUP I (Y) I.
8) cn[/m] is the space of all functions defined in an XEX II x II YEt II y I
m-dimensional cube that are continuously differentiable Even a more general theorem is valid: Let a real-valued
up to and including the order n, with the norm of uni- function p (x) defined on a linear space satisfy the con-
form boundedness in all derivatives of order at most n. ditions:
9) M[a, b] is the space of bounded measurable func- p(x +y) ~ p(x)+p(y),
tions with the norm p(Ax) = 'Ap(x), A;;;' 0, X,YEX,
I x II = ess max I x (t) I and let f (x) be a real-valued linear functional defined
a~t~b
on a subspace Y C X and such that
10) A (D) is the space of functions which are analytic
in the open unit disc D and are continuous in the I(x) ~ p(x), XE Y.
closed disc ]5, with the norm Then there exists a linear functional F(x) defined on
the whole of X such that
II x I = rna]( I x (z) I
::'ED
F(x) = I(x), XE Y; F(x) ~ p(x), XEX
11) Lp(S;~,JL),p;;;;'l, is the space of functions x(s)
A consequence of the Hahn - Banach theorem is the
defined on a set S provided with a countably-additive
'inverse' formula which relates the norms of X and X':
measure JL, with the norm
- I (x) I
IP II x II - ~ 11/11 ,/~O, XEX
IIxll = [fl x(s) IPJL{ds)f .
The maximum in this formula is attained for some
12) Ip[a, b], p;;;;.l, is a special case of the space f= Ix EX'. Another important consequence is the

335
BANACH SPACE

existence of a separating set of continuous linear func- A Banach space is said to be strictly convex if its unit
tionals, meaning that for any x \*x 2 E X there exists a sphere S contains no segments. Convexity moduli are
linear functional f on X such that f(x\)*f(X2) (cf. introduced for a quantitative estimation of the convex-
Complete set of functionals). ity of the unit sphere; these are the local convexity
The general form of a linear functional is known for modulus
many specific Banach spaces. Thus, on Ip[a, b], p> 1,
all linear functionals are given by a formula S(X,f) = inf{l-II X;y II:YES, Ilx-YII;;a.+
b
f(x)= !x(t)y(t)dt,
a XES, 0<f';;;2,
where Y ELq[a, b], 1 / p + 1 / q = 1, and any function and the uniform convexity modulus
y(t)ELq defines a linear functional f by this formula,
moreover b ]I/q 8(f) = inf 8(x, f).
XES

Ilfll = [ [Iy(t) Iqdt . If ~(x, 0 for all x ES and all >0, the Banach space
Thus, the dual space of Lp is Lq: L; = Lq. Linear func-
tionals on L \ [a, b] are defined by the same formula,
but in this case y EM, so that Li =M.
The space X", dual to X', is said to be the second
dual. Third, fourth, etc., dual spaces are defined in a
similar manner. Each element in X may be identified
with some linear functional defined on X':
F(f) = f(x) for allfEX' (FEX", XEX),
where II F II = I x II One may then regard X as a sub-
space of the space X" and X ex" C XIV ... ,
x' ex'" c .... If, as a result of these inclusions, the
Banach space coincides with its second dual, it is called
reflexive. In such a case all inclusions are equalities. If
is said to be locally uniformly convex. If ~(O, the
space is said to be uniformly convex. All uniformly
convex Banach spaces are locally uniformly convex; all
locally uniformly convex Banach spaces are strictly
convex. In finite-dimensional Banach spaces the
converses are also true. If a Banach space is uniformly
convex, it is reflexive.
A Banach space is said to be smooth if for any
linearly independent elements x and y the function
1f;(t) = II x + ty II is differentiable for all values of t. A
Banach space is said to be uniformly smooth if its
modulus of smoothness
per) = sup { I X +ry "; II X -ry II
X.YES
-I}, r>O,

X is not reflexive, all inclusions are strict. If the quo- satisfies the condition
tient space X" / X has finite dimension n, X is said to
be quasi-reflexive of order n. Quasi-reflexive spaces exist
lim 2S2l = o.
for all n. In uniformly smooth spaces, and only in such spaces,
RefleXiVity criteria for Banach spaces. 1) X is reflex- the norm is uniformly Frechet differentiable. A uni-
ive if and only if for each f EX' it is possible to find an formly smooth Banach space is smooth. The converse is
x EX on which the 'sup' in the formula true if the Banach space is finite-dimensional. A
II f II -- ~~~ I I[(x) I
x II ' X=F0,
Banach space X is uniformly convex (uniformly
smooth) if and only if X is uniformly smooth (uni-
is attained. formly convex). The following relationship relates the
2) In reflexive Banach spaces and only in such spaces convexity modulus of a Banach space X and the
each bounded set is relatively compact with respect to smoothness modulus of X':
weak convergence: Anyone of its infinite parts con-
pX'(r) suo {f(2T)
= 0<,l;;2 -8 X(f)}.
tains a weakly convergent sequence (the
Eberletn - Shmul'yan theorem). The spaces Lp and Ip' If a Banach space is uniformly convex (uniformly
p> 1, are reflexive. The spaces L \, 11, C, M, c, m, AP smooth), so are all its subspaces and quotient spaces.
are non-reflexive. The Banach spaces Lp and Ip, p> 1, are uniformly con-
A Banach space is said to be weakly complete if each vex and uniformly smooth, and
weak Cauchy sequence in it weakly converges to an ele-
2 (I <p .;;;2),
ment of the space. Every reflexive space is weakly com- {
8() c:o::: f!' (2.;;;p<oo);
plete. Moreover, the Banach spaces L 1 and 11 are
weakly complete. The Banach spaces not containing a i' (l<p.;;;2),
{
subspace isomorphic to Co form an even wider class. per) c:o::: ~ (2';;;p<oo);
These spaces resemble weakly-complete spaces in
several respects.
lrf()c:o:::CP() ~ a <M<hl.
CP() J

336

The Banach spaces M, C, A, L I , AP, m, c, II are not
strictly convex and are not smooth.
The following important theorems for linear opera-
tors are valid in Banach spaces:
The Banach - Steinhaus theorem. If a family of linear
operators !if = {T a} is bounded at each point,
sup II Tax II < 00, XEX,
a
then it is norm-bounded:
sup II Ta II < 00.
a
The Banach open-mapping theorem. If a linear con-
tinuous operator maps a Banach space X onto a
Banach space Y in a one-to-one correspondence, the
inverse operator T- I is also continuous.
The closed-graph theorem. If a closed linear operator
maps a Banach space X into a Banach space Y, then it
is continuous.
Isometries between Banach spaces occur rarely. The
classical example is given by the Banach spaces L2 and
12 The Banach spaces C[Kd and C[K 2 ] are isometric
if and only if K] and K2 are homeomorphic (the
Banach-Stone theorem). A measure of proximity of
isomorphic Banach spaces is the number
d(X, Y)=lninf II T 11'11 T- 1 II,
T
where T runs through all possible operators which real-
ize a (linear topological) isomorphism between X and
Y. If X is isometric to Y, then d(X, Y)=O. However,
non-isometric spaces for which d(X, Y)=O also exist;
they are said to be almost-isometric. The properties of
Banach spaces preserved under an isomorphism are
said to be linear topological. They include separability,
reflexivity and weak completeness. The isomorphic clas-
sification of Banach spaces contains, in particular, the
following theorems:
Lr =/=- Ls; Ir =/=- I" r=/=-s;
Lr =/=-1" r=/=-s; Lr = I,. r=s=2;
M = m; C[O,I)=/=-A(D);
C[K] = C [0, 1] if K is a metric compactum with the
cardinality of the continuum;
C[Im) =/=- C[O, I).
Each separable Banach space is isomorphic to a locally
uniformly convex Banach space. It is not known (1985)
if there are Banach spaces which are isomorphic to
none of their hyperplanes. There exist Banach spaces
which are not isomorphic to strictly convex spaces.
Irrespective of the linear nature of normed spaces, it is
possible to consider their topological classification. Two
spaces are homeomorphic if a one-to-one continuous
correspondence, such that its inverse is also continuous,
can be established between their elements. An incom-
BANACH SPACE
plete normed space is not homeomorphic to any
Banach space. All infinite-dimensional separable
Banach spaces are homeomorphic.
In the class of separable Banach spaces, C [0, 1] and
A (D) are universal (cf. Universal space). The class of
reflexive separable Banach spaces contains even no iso-
morphic universal spaces. The Banach space II is
universal in a somewhat different sense: All separable
Banach spaces are isometric to one of its quotient
spaces.
Each of the Banach spaces mentioned above, except
L2 and 12 , contains subspaces without a complement.
In particular, in m and M every infinite-dimensional
separable subspace is non-complementable, while in
C[O, 1] all infinite-dimensional reflexive subspaces are
non-complement able. If all subspaces in a Banach
space are complementable, the space is isomorphic to a
Hilbert space. It is not known (1985) whether or not all
Banach spaces are direct sums of some two infinite-
dimensional subspaces. A subspace Y is complement-
able if and only if there exists a projection which maps
X onto Y. The lower bound of the norms of the projec-
tions on Y is called the relative projection constant
,\( Y, X) of the subspace Y in X. Each n-dimensional
subspace of a Banach space is complementable and
'\(Yn' X)';;;; Vn. The absolute projection constant ,\(Y) of
a Banach space Y is
A(Y) = S~A(Y, X),
where X runs through all Banach spaces which contain
Yas a subspace. For any infinite-dimensional separable
Banach space Y one has '\(Y)= 00. Banach spaces for
which '\(Y)';;;;'\<oo form the class fJ'A ('\;;'1). The
class fJ'1 coincides with the class of spaces C(Q) where
Q are extremally-disconnected compacta (d.
Extremally-disconnected space).
Fundamental theorems on finite-dimensional Banach
spaces. 1) A finite-dimensional space (a Minkowski
space) is complete, i.e. is a Banach space. 2) All linear
operators in a finite-dimensional Banach space are con-
tinuous. 3) A finite-dimensional Banach space is reflex-
ive (the dimension of X is equal to the dimension of
X). 4) A Banach space is finite-dimensional if and only
if its unit ball is compact. 5) All n-dimensional Banach
spaces are pairwise isomorphic; their set becomes com-
pact if one introduces the distance
d(X, Y) = lninf II T 11'11 r- I II
T
A series
(*)
is said to be convergent if there exists a limit S of the
sequence of partial sums:
lim
n_oo
Ils- ~>k I = 0.
k =I
337
BANACH SPACE
If
the series (*) is convergent, and is said in such a case to
be absolutely convergent. A series is said to be uncondi-
tionally convergent if it converges when its terms are
arbitrarily rearranged. The sum of an absolutely con-
vergent series is independent of the arrangement of its
terms. In the case of series in a finite-dimensional
space (and, in particular, for series of numbers) uncon-
ditional and absolute convergence are equivalent. In
infinite-dimensional Banach spaces unconditional
convergence follows from absolute convergence but the
converse is not true in any infinite-dimensional Banach
space. This is a consequence of the Dvoretskii- Rogers
theorem: For all numbers ak ;;"0, subject to the condi-
tion ~at < 00, there exists in each infinite-dimensional
Banach space an unconditionally convergent series
~Xk such that I Xk II =ab k = 1, 2, .... In the space
Co (and hence also in any Banach space containing a
subspace isomorphic to co), for any sequence ak;;"O
that converges to zero, there exists an unconditionally
convergent series ~Xb II Xk II =ak' In Lp(S;~, f.L) the
unconditional convergence of the series ~Xk implies
that 00
~ II Xk liS < 00,
k =1
where
s _
-
{2P (l';;;p';;;2),
(P~2).
In a uniformly convex Banach space with convexity
modulus 8(i) the unconditional convergence of the
series ~Xk implies that
00
~o(11 Xk II) < 00.
k =1
A series ~Xk is said to be weakly unconditionally
Cauchy if the series of numbers ~ If (xd I converges
for each fEX. Each weakly unconditiona) y Cauchy
series in X converges if and only if X contains no sub-
space isomorphic to co.
r
A sequence of elements {ek} of a Banach space is
said to be minimal if each one of its terms lies outside
the closure of x(n) = [ekh""'n' the linear hull of the
remaining elements. A sequence is said to be uniformly
minimal if
p(en;x(nJ)~Yllenll, O<y,;;;l, n=1,2, ....
If y = 1, the series is said to be an A uerbach system. In
each n-dimensional Banach space there exists a com-
plete Auerbach system {e n }7. It is not known (1985)
whether or not a complete Auerbach system exists in
each separable Banach space. For each minimal system
there exists an adjoint system of linear functionals {j,,},
which is connected with {ek} by the biorthogonality
33~
relations: Ji(ej)=8ij. In such a case the system {ek,fk}
is said to be biorthogonal. A set of linear functionals is
said to be total if it annihilates only the zero element of
the space. In each seperable Banach space there exists
a complete, minimal system with a total adjoint. Each
element x EX can formally be developed in a series by
the biorthogonal system:
00
x ~ ~.fk(x)eb
k =1
but in the general case this series is divergent.
r
A system of elements {ek} is said to be a basis in X
if each element x EX can be uniquely represented as a
convergent series
00
x = ~ akeb ak =ak(x),
k=1
Each basis in a Banach space is a complete uniform
minimal system with a total adjoint. The converse is
not true, as can be seen from the example of the system
{eint}~oo in crO,2'IT] and LdO,2'IT].
A basis is said to be unconditional if all its rearrange-
ments are also bases; otherwise it is said to be condi-
tional. The system {eint}~oo in Lp [0,2'IT],p>l,p=2, is
a conditional basis. The Haar system is an uncondi-
tional basis in Lp ' p> 1. There is no unconditional
basis in the spaces C and L \. It is not known (1985)
whether or not each Banach space contains an infinite-
dimensional subspace with an unconditional basis. Any
non-reflexive Banach space with an unconditional basis
contains a subspace isomorphic to 1\ or co.
Two normalized bases {e~} and {e~} in two Banach
spaces X \ and X 2 are said to be equivalent if the
correspondence e~e~, k = 1,2, ... , may be extended to
an isomorphism between X \ and X 2. In each of the
spaces 12 , 1\, Co all normalized unconditional bases are
equivalent to the natural basis. Bases constructed in
Banach spaces which have important applications are
not always suitable for solving problems, e.g. in the
theory of operators. T-bases, or summation bases, have
been introduced in this context. Let {tij} r be the
matrix of a regular summation method (d. Regular
summation methods). The system of elements {en} eX
is said to be a T-basis corresponding to the given sum-
mation method if each x EX can be uniquely
represented by a series
oc
X ~ ~akeko
k =I
which is summable to x by this method. The tri-
gonometric system {e 'nl } ~ x in C[O, 27T] is a summation
basis for the methods of Cesaro and Abel. Each T-
basis is a complete minimal (not necessarily uniformly
minimal) system with a total adjoint. The converse is
not true. Until recently (the 1970's) one of the principal

problems of the theory of Banach spaces was the basis
problem dealt with by Banach himself: Does a basis
exist in each separable Banach space? The question of
existence of a basis in specifically defined Banach
spaces remained open as well. The first example of a
separable Banach space without a basis was constructed
in 1972; bases in the spaces cn(Im) and A(D) have
been constructed.
References
[I] BANACH, S.: 'Sur les operations dans les ensembles abstraits et
leur application aux equations integrales', Fund. Math. 3
(1922), 133-181.
[2] BANACH, S.S.: A course offunctionol analysis, Kiev, 1948 (in
Ukrainian).
[3] DUNFORD, N. and SCHWARTZ, J.T.: Linear operators. General
theory, I, Interscience, 1958.
[4] DAY, M.M.: Normed linear spaces, Springer, 1958.
[5] BOURBAKI, N.: Elements of mathematics. Topological vector
spaces, Addison-Wesley, 1977 (translated from the French).
[6] SINGER, I.: Bases in Banach spaces, 12, Springer, 1970-1981.
[7] LINDENSTRAUSS, J. and TZAFRIRI, L.: Classical Banach spaces,
1-2, Springer, 1977-1979.
[8] DIESTEL, J.J.: Geometry of Banach spaces. Selected topics,
Springer, 1975.
[9] BEAUZAMY, B.: Introduction to Banach spaces and their
geometry, North-Holland, 1985.
M.l. Kadets
B.M. Levitan
Editorial comments. The second dual of a space is also
called the bidual.
References
[A 1] SEMANEDl, Z.: Banach spaces of continuous functions, Polish
Sci. Publ, 1971.
AMS 1980 Subject Classification: 468XX, 46EXX
BANACH - STEINHAUS THEOREM - A general
appellation for several results concerning the linear-
topological properties of the space of continuous linear
mappings of one linear topological space into another.
Let E and F be locally convex linear topological spaces,
where E is a barrelled space, or let E and F be linear
topological spaces, where E is a Baire space. The fol-
lowing propositions are then valid. I) Any subset of the
set L(E, F) of continuous linear mappings of E into F
which is bounded in the topology of simple conver-
gence is equicontinuous (the uniform boundedness princi-
ple); 2) If a filter P in L(E, F) contains a set bounded
in the topology of simple convergence, and converges
in the topology of simple convergence to some mapping
v of E into F, then v is a continuous linear mapping of
E into F, and P converges uniformly to v on each com-
pact subset of E [2], [3].
These general results make it possible to render the
classical results of S. Banach and H. Steinhaus [I] more
precise: Let E and F be Banach spaces and let M be a
subset of the second category in E. Then, I) if
HCL(E, F) and sup{ II u(x) II: uEH} is finite for all
xEM. then sup {II u II: uEH}<oc; 2) if Un IS a
BAND OF SEMI-GROUPS
sequence of continuous linear mappings of E into F,
and if the sequence un(x) converges in F for all x EM,
then Un converges uniformly on any compact subset of
E to a continuous linear mapping v of E into F.
References
[I] BANACH, S. and STEINHAUS, H.: 'Sur Ie principe de la conden-
sation de singularites', Fund. Math. 9 (1927),50-61.
[2] BOURBAKI, N.: Elements of mathematics. Topological vector
spaces, Addison-Wesley, 1977 (translated from the French).
[3] SCHAEFER, H.H.: Topological vector spaces, Macmillan, 1966.
A.I. Shtern
Editorial comments.
References
[A 1] KOTHE, G.: Topological vector spaces, 1, Springer, 1969.
[A2] KELLEY, 1.L. and NAMIOKA, I.: Linear topological spaces,
Springer, 1963.
AMS 1980 Subject Classification: 46A32
BAND OF SEMI-GROUPS of a given family {Sa} -
A semi-group S that has a partition into sub-semi-
groups whose (isomorphism) classes are just the semi-
groups Sa, and such that for any Sa, S f3 there is an S y
such that Sa S f3 CS y' S is also said to be decomposable
into the band of semi-groups Sa. In other words, S has
a partition into a band of semi-groups Sa if all the Sa
are sub-semi-groups of S and if there is a congruence p
on S such that the p-classes are just the Sa. The semi-
groups Sa are called the components of the given band.
The term 'band of semi-groups' is consistent with the
frequent use of the word 'band' as a synonym of 'semi-
group all elements of which are idempotents', since a
congruence p on a semi-group S determines a partition
of S into a band if and only if the quotient semi-group
S / p is a semi-group of idempotents.
Many semi-groups are decomposable into a band of
semi-groups with one or other 'better' property; thus,
the study of their structure is reduced in some measure
to a consideration of the types to which the com-
ponents of a band belong, and of semi-groups of idem-
potents (see, e.g. Archimedean semi-group; Completely-
simple semi-group; Oifford semi-group; Periodic semi-
group; Separable semi-group).
A band of semi-groups Sa is said to be commutative
if for the corresponding congruence p the quotient
semi-group S / p is commutative; then S / p is a semi-
lattice (in this case, S is frequently called a semi-lattice
of semi-groups Sa; in particular, if S / p is a chain,
then S is called a chain of semi-groups Sa). A band of
semi-groups is called rectangular (sometimes matrix) if
S / p is a rectangular semi-group (see Idempotents,
semi-group of). Equivalently, if the components of the
band can be indexed by pairs of indices Si"A, where i
and A run over certain sets I and A, respectively, such
that for any So" Sj/" one has Si"ASJ/" CSi/". Any band of
semi-groups is a semi-lattice of rectangular bands, that
339
BAND OF SEMI-GROUPS
is, its components can be arranged into subfamilies so
that the union of the components of each subfamily is
a rectangular band of components, and the original
semi-group is decomposable into a semi-lattice of these
unions (Clifford's theorem [1]). Since the properties of
being a semi-group of idempotents, a semi-lattice or a
rectangular semi-group are characterized by identities,
for each of the listed properties () there is a finest
congruence on any semi-group S for which the
corresponding quotient semi-group has the property (),
that is, there exist greatest (or biggest quotient) partitions
of S into a band of semi-groups, into a commutative
band of semi-groups and into a rectangular band of
semi-groups.
The term strong band concerns special types of bands
of semi-groups [4]: For any elements a and b from
different components, the product ab is a power of one
of these elements. An important special case of a
strong band, and also a special case of a chain of
semi-groups, is the ordinal sum (or sequentially-
annihilating band): The set of its components {Sa} is
totally ordered, and for any Sa, S p such that Sa <S p,
and for any aES a , bESp one has ab=ba=a; the ordi-
nal sum is defined uniquely up to an isomorphism, by
specifying the components and their ordering.
References
[I] CLIFFORD, A.H.: 'Bands of semi-groups', Proc. A mer. Math.
Soc. 5 (1954), 499-504.
[2] CLIFFORD, A.H. and PRESTON, G.B.: The algebraic theory of
semigroups, I, Amer. Math. Soc., 1961.
[3] LYAPIN, E.S.: Semigroups, Amer. Math. Soc., 1974 (translated
from the Russian).
[4] SHEVRIN, L.N.: 'Strong bands of semi-groups', lzv. Vyssh.
Uchebn. Zaved. Mat., no. 6 (1965), 156-165 (in Russian).
L.N. Shevrin
Editorial comments. A congruence on a semi-group 5 is
an equivalence relation such that for all a, b, CES one has
a~b =} ac~bc and ca~cb.
AMS 1980 Subject Classification: 20M10
BAR INDUCTION - An inductive way of reasoning,
used in intuitionistic mathematics (cf. Intuitionism),
which may be described as follows. Let certain proper-
ties Rand S be defined on finite tuples (cf. Tuple) of
natural integers, such that I) the property R is decid-
able, i.e. it is possible to effectively find out whether it
is fulfilled on any tuple <n 1, . . . ,nm > or not; 2) for
any choice sequence 0: it is possible to find a tuple of
the form <0:(0), ... ,o:(n -I for which R is fulfilled.
Moreover, if 2) is fulfilled, one says that R 'bars' the
empty tuple < > (hence the name); 3) for any tuple 1T
of natural integers, if R(1T), then S(1T) (the so-called
basis of the bar induction); and 4) if <n], ... ,n m > is
a tuple such that for any natural number k one has
340
SnJ, ... ,nm , k, then Sn), ... ,nm ; this
property is known as the step of the bar induction.
If the conditions 1) through 4) above are met, one
can deduce from the principle of bar induction that
S .
The use of bar induction was proposed by L.E.1.
Brouwer as an intuitionistic method of reasoning which
indicates that the collection of free choice sequences is
incomplete and to some extent uncountable. It was
demonstrated, in particular, by S.c. Kleene and,
independently, by AA Markov that it follows from the
principle of bar induction (in fact, even from a corol-
lary of bar induction, the fan theorem) that not all
choice sequences are recursive (d. Fan).
The forms of bar induction which one started to util-
ize in the foundations of mathematics in the 1960's do
not deal with tuples of natural numbers, but rather
with tuples of more complex objects, such as tuples of
choice sequences.
Bar induction can be written down in the language
of formal intuitionistic mathematical analysis as fol-
lows:
'Vx (R(x)V -,R(x & 'Va 3x R(a(x &
& 'Vx (R(xPS(x &
& 'Vx ('Vy S(x*'yPS(x:J S(Q).
References
[I] KLEENE, S.c. and VESLEY, R.E.: Thefoundations of intuition is-
tic mathematics: especially in relation to recursive functions,
North-Holland, 1965.
A. G. Dragalin
Editorial comments.
References
[A 1] DUMMETJ, M.: Elements of intuitionism, Clarendon Press,
1977.
[A2] TROELSTRA, A.S.: Choice sequences, Clarendon Press, 1977.
AMS 1980 Subject Classification: 03F55
BARBER PARADOX - The same as the antinomy of
the 'village barber'.
AMS 1980 Subject Classification: 03A05, 00A25
BARBIER THEOREM on curves of constant width - If
the distance between any two parallel supporting
straight lines to a curve is constant and equal to a. then
the length of the curve is 1Ta. Discovered by E. Barbier
in 1860.
A.B. Ivanov
Editorial comments. The original work of E. Barbier IS
[A 1].
References
[A1] BARBIER. E.: 'Note sur Ie probleme de l'amgullle et Ie leu du
Joint couvert', J Math. Pure Appl 5 (1860), 273-286.
[A2] BO""lsF~. T. and FEVHH. W. Theorie der konvexen
K6rper, Springer. 1934.
AMS 1980 Subject Classification: 53A04, 52A 10

BARRELLED SPACE - A locally convex linear topo-
logical space displaying several properties of Banach
spaces and Frechet spaces without the metrizability
condition (cf. Banach space; Frechet space). It is one of
the most extensive class of spaces to which the
Banach - Steinhaus theorem applies. Barrelled spaces
were first introduced by N. Bourbaki.
A set A in a vector space E is said to be a balanced
set if ax EA for all x EA and for all a such that
I a I ~ 1. A balanced set ACE is said to be an absorb-
ing set if it absorbs each point of E, i.e. if for each
x EE there exists an a>O such that ax EA.
A barrel in a linear topological space is a closed, bal-
anced, absorbing, convex set. A barrelled space is a
linear topological space with a locally convex topology
in which every barrel is a neighbourhood of zero.
Frechet spaces and, in particular, Banach spaces are
examples of barrelled spaces. Montel spaces (cf. Montel
space) are an important class of barrelled spaces, and
display remarkable properties. A quotient space of a
barrelled space, a direct sum and inductive limits of
barrelled spaces are barrelled spaces. Every pointwise-
bounded set of continuous linear mappings of a bar-
relled space into a locally convex linear topological
space is equicontinuous. In a space dual to a barrelled
space, a bounded set in the weak topology is bounded
in the strong topology and relatively compact in the
weak topology.
References
[1] BOURBAKI, N.: Elements of mathematics. Topological vector
spaces, Addison-Wesley, 1977 (translated from the French).
[2] EDWARDS, R.E.: Functional analysis: theory and applications.,
Holt, Rinehardt, Winston, 1965.
V. M. Tikhomirov
Editorial comments. Barrelled spaces are the most
extensive class of locally convex spaces to which the
Banach-Steinhaus theorem can be extended. They were
first introduced in [A4).
A not necessarily balanced set A in E is called absorbing
if for every xEE there is an ao such that xEaA for all
I a I ;;"aa For the dual of a barrelled space the following
four statements are equivalent: 1) A is weakly bounded; 2)
A is strongly bounded; 3) A is equi-continuous; and 4) A is
weakly compact. The last statement follows from the
stronger statement that the dual of a barrelled space is
quasi-complete for any a-topology. (For the last notion see
Topological vector space; Space of mappings, topologi-
cal.)
References
[A 1] ScHAEFER, H.H.: Topological vector spaces, Macmillan,
1966.
[A2] KELLEY, J.L. and NAMIOKA, I.: Linear topological spaces,
Springer, 1963.
[A3] KOTHE, G.: Topological vector spaces, 1, Springer, 1969.
[A4] BoURBAKI, N.: 'Sur certains espaces vectoriels topologiques',
Ann. Inst. Fourier 2 (1950), 5-16.
AMS 1980 Subject Classification: 46A07
BARRIER
BARllIER, Lebesgue barrier, in potential theory - A
function the existence of which is a necessary and suffi-
cient condition for the regularity of a boundary point
with respect to the behaviour of a generalized solution
of the Dirichlet problem at that point (cf. Perron
method; Regular boundary point).
Let D be a domain in a Euclidean space Rn, n ~ 2,
and let ~ be a point on its boundary f=aD. A barrier
for the point ~ is any function w(x), continuous in the
intersection (D U f) n
B of the closed domain D U f
with some ball B =B(R, ~ with centre at ~, which is
superharmonic in D n
B and positive in (D U f) B, n
except at ~, at which it vanishes. For instance, if n ~3
and ~ is any boundary point for which there exists a
closed ball B(R, y) in D U f which meets f only in ~,
one can take as a barrier the harmonic function
1 1
w~(x) = Rn-2 - I X -y In- 2'
where R is the radius of B(R, y) and y is its centre.
A barrier in the theory of functions of (several) com-
plex variables is a function the existence of which for all
boundary points of the domain D implies that D is a
domain of holomorphy. Let D be a domain in the com-
plex space en, n ~ 1, and let ~ be a point of the boun-
dary f=aD. Any analytic function fez) in D with a
singular point at ~ will then be a barrier at ~. Thus, the
function 1 / (z -~) is a barrier for the boundary point ~
of any plane domain DeC. There also exists a barrier
at any point ~=(~l' ... '~n) of the boundary of the ball
= {z = (Zh ,zn): I
D Zi 1 2 <R 2 },
I~l
e.g. the function I / (~l Z I + ... +Inzn - R 2).
A barrier exists at a boundary point ~ of a domain D
if there is an analytic function defined in D that is
unbounded at ~, i.e. is such that for some sequence of
points {z(k)} CD which converges to ~ one has:
lim
k ..... oo
If (Z(k) I = + 00.
The converse is true for domains D C en in the fol-
lowing strong form: For any set E of boundary points
of a domain D at which a barrier exists, one can find a
function holomorphic in D which is unbounded at all
points of E. If E is everywhere dense in the boundary
of D, then D is a domain of holomorphy.
References
[1] COURANT, R. and HILBERT, D.: Methods of mathematical phy-
sics. Partial differential equations, 2, Interscience, 1962
(translated from the German).
[2] VLADIMIROV, V.S.: Methods of the theory offunctions of several
complex variables, M.I.T., 1966, Chapt. 3 (translated from
the Russian).
[3] SHABAT, B.Y.: Introduction to complex analysis, 1-2, Moscow,
1985, Chapt. 3 (in Russian).
E.D. Solomentsev
M. Shirinbekov
341
BARRIER
Editorial comments. Good English references for the
Lebesgue barrier are [A1] and [A2].
References
[A1] HAYMAN, W.K. and KENNEDY, P.B.: Subharmonic functions,
1, Cambridge Univ. Press, 1976.
[A2] HELMS, L.L.: Introduction to potential theory, Acad. Press,
1975 (translated from the German).
AMS 1980 Subject Classification: 31825, 31 A25,
32A40
BARTLETf TEST - A statistical test in the
Behrens- Fisher problem. It refers to the case of two
samples Xl, . . . ,Xn , EN(a\, aT) and
y\, ... 'Yn, EN(a2' a~) having equal size, n\ =n2 =n.
The critical regions are given by the inequality
Vn i--
----~--~~~===-------~ > c
for some constant c >0, where
_ 1 n _ 1 n
x = - ~Xi' Y = - ~Yi'
n i=1 n 1=1
The left-hand side of the inequality has a Student
distribution with n - 1 degrees of freedom, connecting
the significance level of the test to the constant c.
Bartlett's test is a special case of the Scheffe test and
has analogous extremal properties, but was presented
by M.S. Bartlett [1] prior to this solution. A second test
of Bartlett is connected with the comparison of the
variances of many samples.
References
(1) BARTLETT, M.S.: 'The information available in small samples',
Proc. Cambridge Phi/os. Soc. (1) 32, no. 3 (1936), 560-566.
0. V. Shalaevskif
AMS 1980 Subject Classification: 62FXX
BARYCENTRIC COORDINATES - CoordiT1ates of a
point in an n-dimensional vector space En, v'th respect
to some fixed system Po, ... ,pn of points that do not
lie in an (n -I )-dimensional subspace. Every point
x EE n can uniquely be written as
x = I\)PII + ... +t..II PII'
where Ao, ... ,An are real numbers satisfying the condi-
tion Ao + ... + An = I. The point x is by definition the
centre of gravity of the masses Ao, ... , An located at
the points po . ... ,pn' The numbers Ao, ... ,An are
called the barl'centric coordinates of the point x; the
point with barycentric coordinates I / (n + I) is called
the barvcentre. Barycentric coordinates were introduced
by A.F. Mobius in 1827. [I]. as an answer to the ques-
tion about the masses to be placed at the vertices of a
triangle so that a given point is the centre of gravity of
these masses. Barycentric coordinates are a special case
342
of homogeneous coordinates; they are affine invariants.
Barycentric coordinates of a simplex are used in alge-
braic topology [2]. Barycentric coordinates of a point of
an n-dimensional simplex a with respect to its vertices
A a, ... ,An is the name given to its (ordinary) Carte-
sian coordinates III the basis of the vectors
OAa, ... ,0A n , where 0 is any point that does not lie
in the n-dimensional subspace carrying a (if it is con-
sidered that a lies in some Euclidean space, then the
definition does not depend on the point 0), or to pro-
jective coordinates with respect to A a, ... ,An in the
projective completion of the subspace containing a. The
barycentric coordinates of the points of a simplex are
non-negative and their sum is equal to one. If the i-th
barycentric coordinate becomes zero, this means that
the point lies at the side of the simplex a opposite to
the vertex Ai' This makes it possible to consider the
barycentric coordinates of the points of a geometric
complex with respect to all of its vertices. Barycentric
coordinates are used to construct the barycentric subdi-
vision of a complex.
Barycentric coordinates of abstract complexes are
formally defined in an analogous manner [3].
References
[I] MOBIUS. A.F.: 'Der barycentrische Kalkul', in Gesammelte
Werke, Vol. I, Hirzel. Leipzig, \885.
[2] PONTRYAGIN. L.S.: Grundzuge der kombinatorischen Topologie,
Deutsch. Verlag Wissenschaft., 1956 (translated from the Rus-
sian).
[3] SPANIER, E.H.: Algebraic topology, MacGraw-Hill, 1966.
E. G. Sklyarenko
AMS 1980 Subject Classification: 51 N1 0, 51 N20
BARYCENTRIC SUBDIVISION of a geometric com-
plex K - A complex K I obtained by replacing the sim-
plices of K by smaller ones by means of the following
procedure. Each one-dimensional simplex (segment) is
halved. On the assumption that all simplices of dimen-
sion ~n - I are already subdivided. the subdivision of
any n-dimensional simplex a is defined by means of
cones over the simplices of the boundary of a with a
common vertex that is the barycentre of the simplex a,
i.e. the point with barycentric coordinates 1/ (n + I).
The vertices of the resulting complex K I are in a one-
to-one correspondence with the simplices of the com-
plex K, while the simplices of the complex K I are in
such a correspondence with inclusion-ordered finite
tuples of simplices from K. The formal definition of a
barycentric subdivision for the case of an abstract com-
plex is analogous.
E. C. Sklvarenko
Editorial comments.
References
[A1 j LA!-.IOTKI, K. SemislmpilZla/e aigebralsche Top%gle,
Springer, 1968
AMS 1980 Subject Classification 51 N 10, 51 N20

BASE of a topological space X, base of a topology,
basis of a topology, open base - A family IB of open sub-
sets of X such that each open subset G C X is a union
of sub collections U C lB. The concept of a base is a fun-
damental concept in topology: in many problems con-
cerned with open sets of some space it is sufficient to
restrict the considerations to its base. A space can have
many bases, the largest one of which is the family of all
open sets. The minimum of the cardinalities of all bases
is called the weight of the topological space X. In a
space of weight 7" there exists an everywhere-dense set
of cardinality ';;;;7". Spaces with a countable base are
also referred to as spaces satisfYing the second axiom of
countability. The dual concept of a closed base, formed
by the complements of the elements of a base, is not
used to any significant extent.
A local base of a space X at a point x EX (a base of
the point x) is a family IB(x) of open sets of X with the
following property: For any neighbourhood Ox of x it
is possible to find an element V E IB(x) such that
x EVe Ox' Spaces with a countable local base at every
point are also referred to as spaces satisfYing the first
axiom of countability. A family IB of open sets in X is a
base if and only if it is a local base of each one of its
points x EX.
Let m, n be cardinal numbers. A base IB of the space
X is called an m-point base if each point x EX belongs
to at most m elements of the family IB; in particular, if
m = 1, the base is called disjoint; if m is finite, it is
called bounded point finite; and if m = ~o, it is called
point countable.
A base IB of the space X is called m-local if each
point x EX has a neighbourhood Ox intersecting with
at most m elements of the family IB; in particular, if
m = 1, the base is referred to as discrete; if m is finite, it
is called bounded locally finite; and if m = ~o, it is called
locally countable. A base IB is called an (n - m)-point
base (or an (n-m)-local base) if it is a union of a set of
cardinality n of m-point (m-Iocal) bases; examples are,
for n=~o, a-disjoint, a-point finite, a-discrete and a-
locally finite bases.
These concepts are used mainly in the criteria of
metrizable spaces. Thus, a regular space with a count-
able base, or satisfying the first axiom of countability
and with a point countable base, is metrizable; a regu-
lar space with a a-discrete or a-locally finite base is
metrizable (the converse proposition is true in the
former case only).
A base IB of the space X is called uniform (k-uniform)
if for each point x EX (each compact subset F) and for
each one of the neighbourhoods Ox (OF) only a finite
number of elements of the base contain x (intersect
with F) and at the same time intersect with the
complement X\ Ox (X\ OF)' A space X is metrizable
BASE CHANGE
if and only if it is paracompact with a uniform base (a
Kolmogorovor To-space with a k-uniform base).
A base IB of the space X is called regular if for each
point x EX and an arbitrary neighbourhood Ox of it
there exists a neighbourhood O~ such that the set of all
the elements of the base which intersect both with O~
and X \ Ox is finite. An accessible or T I-space is
metrizable if and only if it has a regular base.
A generalization of the concept of a base is the so-
called w-base (lattice base), which is a family IB of open
sets in the space X such that each non-empty open set
in X contains a non-empty set from IB, i.e. IB is dense
in X according to Hausdorff. All bases are w-bases, but
the converse is not true; thus, the set Z+ in the
Stone- Cech compactification of the set of natural
numbers in Z+ forms only a w-base.
References
[I] ALEKSANDROV, P.S.: EinfUhrung in die Mengenlehre und die
Theorie der reelen Funktionen, Deutsch. Verlag Wissenschaft.,
1956 (translated from the Russian).
[2] URYSOHN, P.S. [P.S. URYSON]: Works on topology and other
fields of mathematics, Vol. 1-2, Leningrad, 1951 (in Russian).
[3] ALEKSANDROV, P.S. and PASYNKOV, B.A.: An introduction to
the theory of topological spaces and general dimension theory,
Moscow, 1973 (in Russian).
[4] ARKHANGEL'SKll, A.V. and PONOMAREV, V.I.: Fundamentals of
general topology: problems and exercises, Reidel, 1984
(translated from the Russian).
[5] BOURBAKI, N.: Elements of mathematics. General topology,
Addison-Wesley, 1966 (translated from the French).
A.A. Mal'tsev
Editorial comments. Closed bases are useful in com-
pactification theory, cf. Compactification.
Besides the notions of a bounded point-finite base and a
bounded local-finite base one also uses that of a point-finite
base and a local-finite base. A base (or any family of sub-
sets ~) is called point finite if every point x belongs to fin-
itely many members of ~, i.e. if ~x={BE~: xEB} is finite
for every x. Note that the families ~x can have arbitrary
large finite cardinalities, in contrast to the definition of
bounded point finiteness, when the cardinalities of ~x are
bounded by a fixed finite m. Similar remarks apply to local
finiteness.
AMS 1980 Subject Classification: 54AXX
BASE CHANGE, change of base - A category-
theoretical construction; special cases are the concept
of an induced fibration in topology, and extension of
the ring of scalars in the theory of modules.
Let C be a category with fib red products and let
g: S I ---"S be a morphism of C. A base change by means
of g is a functor from the category of S-objects (i.e. the
category of morphism f: X ---"S, where X is an object of
C) to the category of S I -objects, taking an S-object
f: X ---"S to the S I-object fl : X I ---"S I, where
X I = X X sS I and fl is projection onto the second fac-
tor. The morphism g is then called the base-change mor-
343
BASE CHANGE
phism. One also says that X I is obtained from X by
base change.
A special case of a base change is the concept of a
fibre of a morphism f: X ~S of schemes S: The fibre of
the morphism f over a point s E S is the scheme
Xs = xXss,
i.e. the scheme obtained from X by base change via the
natural morphism S~S. A similar definition yields the
geometric fibre Xs; it is obtained by base change via the
morphism Spec K ~S associated with a geometric point
of S, where K is an algebraically closed field. Many
properties of the S-scheme X are preserved under a
base change. The inverse problem - to infer the pro-
perties of a scheme X from those of the schemes
obtained from X by base change - is considered in
descent theory (see also [3]).
Let fl: X I ~S I be the morphism obtained from
f: x~S via a morphism g: SI ~S, so that one has a
Cartesian square g,
XI ~X
!I~ U
SI ~S
)(
Let F be a sheaf of sets on X. Then there exists a
natural sheaf mapping tfi:gj.(F)~fl.gI(F). If Fis a
sheaf of Abelian groups, then for every q ~o there
exists a natural sheaf homomorphism
>fq:g'(Rqf.(F)) ~ Rq!I,(g'I(F)).
Under these conditions, tfi and tfiq are also called base-
change morphisms. It is usually said that the base-
change theorem is valid if tfi (or tfiq) is an isomorphism.
In other words, the base-change theorem is a proposi-
tion about the compatibility (commutability) of the
functors Rqj. with the base-change functor. In particu-
lar, if g is the imbedding of a point s ES, the base-
change theorem states that there exists a natural iso-
morphism (Rqf.(F))x :::::=.Hq(Xx, F IxJ bel),'een the
fibre of the q-th direct image of the sheaf F and the q-
dimensional cohomology group of the fibre of the mor-
phism f The base-change theorem is valid in the
following situations: 1) f is a proper mapping of
paracompact topological spaces, S is a locally compact
space [1]; 2) f is a separable quasi-compact morphism
of schemes, g is a flat morphism, F is a quasi-coherent
sheaf of Ox-modules (the comparison theorem for the
cohomology of ordinary and formal schemes - see [2]
- can also be interpreted as a base-change theorem);
or 3) f is a proper morphism of schemes, F is a torsion
sheaf in the etale topology. Some other cases in which
base-change theorems are valid are considered in [3].
References
[I] GODEMENT, R.: Topologie algebrique et theorie des faisceaux,
Hermann, 1958.
344
[2] GROTHENDIECK, A and DIEUDONN~, J.: 'Elements de
geometrie algebrique', Publ. Math. IHES 11 (1961).
[3] ARTIN, M., GROTHENDIECK. A and VERDIER, J.L. (EDS.):
'Theorie des topos et cohomologie etale des scMmas', in Sem
Geom Alg., Vol. 4, Springer, 1973.
v.l. Danilov
AMS 1980 Subject Classification: 18-XX
BASE OF A DEFORMATION - A conjugate net on a
surface F and its deformation r
outside their points of
congruence. The base of a deformation is characterized
by the fact that the bend - the relation between the
normal curvatures k and k* at isometrically-
corresponding points of F and r along corresponding
directions - has extremal values along the directions
of the base of the deformation.
References
[I] KAGAN, V.F.: Foundations of the theory of surfaces in a tensor
setting, Moscow-Leningrad, 1948, Chapt. 2 (in Russian).
M.l. Voftsekhovskif
Editorial comments. For more references on the topic of
deforming or bending surfaces, cf. the article Deformation,
isometric.
AMS 1980 Subject Classification: 53A07
BASIC COMMUTATOR, regular commutator - An
object inductively constructed from the elements of a
given set R and from brackets, in the following manner.
The elements of R are considered by definition to be
basic commutators of length 1, and they are given an
arbitrary total order. The basic commutators of length
n, where where n > 1 is an integer, are defined and
ordered as follows. If a, b are basic commutators of
lengths smaller than n, then [ab] is considered to be a
basic commutator of length n if and only if the follow-
ing conditions are met: 1) a, b are basic commutators
of lengths k and I, respectively, and k +I=n; 2) a>b;
and 3) if a =[cd], then d~b. The basic commutators of
length not exceeding n thus obtained are arbitrarily
ordered, subject to the condition that lab ]>b, while
preserving the order of the basic commutators of
lengths less than n. Any set of basic commutators con-
structed in this way is a base of the free Lie algebra
with R as set of free generators [1].
References
[I] SHIRSHOY, AI.: 'On bases of free Lie algebras', Algebra i
Logika 1, no. I (1962), 14-19 (in Russian)v M G h k
Lorc a ov
U. .
Editorial comments. Let M(R) be the free magma on R,
i.e. the set of all non-commutative and non-associative
words in the alphabet R. The basic commutators are to be
seen as a subset of M(R). This subset is also often called a
P. Hall set. The identity on R induces a mapping
</>: M(R)-,>LK(R), where LK(R) is the free Lie algebra on R
over the ring K. Let K(R) be a P. Hall set in M(R) (i.e. a set
of basic commutators), then <P(H(R)) is a basis of the free

K-module LK(R), called a P. Hall basis. Other useful bases
of LK(R) are the Chen- Fox- Lyndon basis and the
Shirshov basis (these two are essentially the same), and the
Spitzer- Foata basis; cf. [A4] for these. Let R be finite of
cardinality r= #R. Let I,(n) be the number of basic commu-
tators on R of length n. Then
I,(n) = .l ~p.(d)rn/d
n dl n
where p.: {1, 2, ... }~{ -1,0, 1} is the Mobius function,
defined by p.(1) = 1, p.(k) = 0 if k is divisible by a square, and
p.(p, ... Pm) =(-1)m if p" ... , Pm are distinct prime
numbers.
References
[A1] BOURBAKI, N.: Groupes etalgebres de Lie, Hermann, 1972,
Chapt. 2; 3.
[A2] HALL, M., JR.: The theory of groups, Macmillan, 1959.
[A3] MAGNUS, W., KARRAss, A. and SOLITAR, D.: Combinatorial
group theory, Interscience, 1966.
[A4] VlENNOT, G.: Algebres de Lie libres et monoides libres, Lect.
Notes in Math., 691, Springer, 1978.
AMS 1980 Subject Classification: 17BXX, 08B20
BASIC SET of a linear system - The set of points of
an algebraic variety (or of a scheme) X which belong to
all the divisors of the movable part of the given linear
system L on X.
Example. Let
AoFn(xo,Xj,X2)+AIGn(XO,Xj,X2) = 0
be a pencil of n-th order curves on the projective plane.
The basic set of this pencil then consists of the set of
common zeros of the forms F' and G', where
F'H = Fn, G'H = Gn,
and H is the greatest common divisor of the forms Fn
and Gn
U cf>L: X~P(L) is the rational mapping defined by L,
then the basic set of L is the set of points of indeter-
minacy of cf>L. A basic set has the structure of a closed
sub scheme B in X, defined as the intersection of all
divisors of the movable part of the linear system. The
removal of the points of indeterminacy of cf>L can be
reduced to the trivialization of the coherent sheaf of
ideals defining the sub scheme B (cf. Birational
geometry).
For any linear system without fixed components L
on a smooth projective surface F there exists an integer
no such that if n >no, then the basic set of the com-
plete linear system I nL I is empty (Zariski's theorem).
This is not true in the multi-dimensional case.
References
[1) Algebraic surfaces, Moscow, 1965 (in Russian).
[2) HARTSHORNE, R.: Algebraic geometry, Springer, 1977.
V.A. Iskovskikh
AMS 1980 Subject Classification: 14C20
BASIS of a set X - A minimal subset B that gen-
BASIS
erates it. Generation here means that by application of
operations of a certain class 0 to elements b EB it is
possible to obtain any element x EX. This concept is
related to the concept of dependence: By means of
operations from 0 the elements of X become dependent
on the elements of B. Minimality means that no proper
subset BleB generates X. In a certain sense this pro-
perty causes the elements of B to be independent:
None of the elements b EB is generated by the other
elements of B. For instance, the set of all natural
numbers Zo has the unique element 0 as basis and is
generated from it by the operation of immediate succes-
sion and its iteration. The set of all natural numbers
> 1 is generated by the operation of multiplication
from the basis consisting of all prime numbers. A basis
of the algebra of quaternions consists of the four ele-
ments {I, i, j, k} if the generating operations consist of
addition and of multiplication by real numbers; if, in
addition to these operations, one also includes multipli-
cation of quaternions, the basis will consist of three ele-
ments only - {l, i, j} (because k = ij).
A basis of the natural numbers of order k is a subse-
quence 0 of the set Zo of natural numbers including 0,
which, as a result of k-fold addition to itself (the gen-
erating operation) yields all of Zo. This means that any
natural number n can be represented in the form
where ai EO. For example, every natural number is a
sum of four squares of natural numbers (Lagrange's
theorem), i.e. the sequence of squares is a basis of Zo of
order 4. In general, the sequence of m-th powers of
natural numbers is a basis of Zo (Hilbert's theorem), the
order of which has been estimated by the Vinogradov
method. The concept of a basis of Zo has been general-
ized to the case of arbitrary sequences of numbers, i.e.
functions on Zo.
A set X always contains a generating set (in the
trivial case: X generates X), but minimality may prove
to be principally impossible (such a situation is typical
of classes 0 containing infinite-place operations, in par-
ticular in topological structures, lattices, etc.). For this
reason the minimality condition is replaced by a
weaker requirement: A basis is a generating set of
minimal cardinality. In this context a basis B is defined
as a parametrized set (or population), i.e. as a function
b(t) on a set of indices T with values in X, such that
b(T)=B; the cardinality of T is sometimes called as
the dimension (or rank) of the basis of X. For example,
a countable everywhere-dense set B in a separable
topological space P may be considered as a basis for it;
P is generated from B by the closure operation (which,
incidentally, is related to generation in more general
cases as well, see below).
345
BASIS
A basis for a topology of a topological space X (a
base) is a basis IS of the set of all open subsets in X;
the generation is effected by taking unions of elements
of IS.
A basis of a Boolean algebra m: (a dual base of m: in
the sense of Tarski) is a dense set S (of minimal cardi-
nality) in m:; the generation of m: from S (and hence S
itself) is determined by the condition s~a = V (which
is equivalent to s Ca), where s ES, a Em:, V is the unit
of m: and '~' is the operation of implication. One also
introduces in an analogous manner a basis for a filter
v as a set S such that for an arbitrary a E v there
exists an SES with s Ca.
More special cases of bases of a set X are introduced
according to the following procedure. Let B (X) be the
Boolean algebra of X, i.e. the set of all its subsets. A
generating operator (or a closure operator) J is a map-
ping of B(X) into itself such that if A CB, then
J(A)CJ(B); A CJ(A); JJ(A)=J(A).
An element x E X is generated by a set A if x EJ (A);
in particular, A generates X if J (A) = X. A minimal set
B possessing this property is said to be a basis of X
defined by the operator J. A generating operator J is of
finite type if, for arbitrary A cX and x CX, it follows
from x EJ (A) that x EJ (A 0) for a certain finite subset
A 0 CA; a generating operator J has the property of sub-
stitution if, for any y, Z EX and A C X, both Y fl.J (A)
and y EJ (A U {z }) imply that Z EJ (A U (y }). A gen-
erating operator J of finite type with the substitution
property defines a dependence relation on X, i.e. a sub-
division of B (X) into two classes - dependent and
independent sets; a set A is said to be dependent if
y EJ (A \y) for some YEA, and is said to be indepen-
dent if y fl.J(A \y) for any y EA. Therefore, A is depen-
dent (independent) if and only if some (arbitrary) non-
empty finite subset(s) A 0 CA is dependent (are
independent).
For a set B to be a basis of the set X it) necessary
and sufficient for B to be an independent generating
set for X, or else, a maximal independent set in X.
If A is an arbitrary independent set, and C is an X-
generating set containing A, then there exists a basis B
in X such that A C B C C. In particular, X always has a
basis, and any two bases of it have the same cardinal-
ity.
In algebraic systems X an important role is played
by the concept of the so-called free basis B, which is
characterized by the following property: Any mapping
of B C X into any algebraic system Y (of the same sig-
nature) may be extended to a (unique)
(homo)morphism from X into Y or, which is the same
thing, for any (homo )morphism B: X ~ Y and any set
A C X, the generating operators J x and J y satisfy the
condition:
346
O{1x(A)} = Jy(O{A}).
An algebraic system with a free basis is said to be free.
A typical example is a basis of a (unitary) module M
over a ring K, that is, a free family of elements from M
generating M [3]. Here, a family A = { at: t E T} of ele-
ments of a K-module M is said to be free if ~~tat = 0
(where ~t =0 for all except a finite number of indices t)
implies that ~t =0 for all t, and the generation is real-
ized by representing the elements x as linear combina-
tions of the elements at: There exists a set (dependent
on x) of elements ~t EK such that ~t =0 for all except a
finite number of indices t, and such that the decompo-
sition
x = ~~tat
is valid (i.e. X is the linear envelope of A). In this sense,
the basis M is free basis; the converse proposition is
also true. Thus, the set of periods of a doubly-periodic
function f of one complex variable, which is a discrete
Abelian group (and hence a module over the ring Z),
has a free baSis, called the period basis of f; it consists
of two so-called primitive periods. A period basis of an
Abelian function of several complex variables is defined
in a similar manner.
If K is a skew-field, all bases (in the previous sense)
are free. On the contrary, there exist modules without
a free basis; these include, for example, the non-
principal ideals in an integral domain K, considered as
a K-module.
A basis of a vector space X over a field K is a (free)
basis of the unitary module which underlies X. In a
similar manner, a basis of an algebra A over a field K is
a basis of the vector space underlying A. All bases of a
given vector space X have the same cardinality, which
is equal to the cardinality of T; the latter is called the
algebraic dimension of X. Each element x EX can be
represented as a linear combination of basis elements in
a unique way. The elements ~t(x)EK, which are linear
functionals on X, are called the components (coordi-
nates) of x in the given basis {ad.
A set A is a basis in X if and only if A is a maximal
(with respect to inclusion) free set in X.
The mapping
where ~X<t) = ~t(x) if ~I is the value of the t-th com-
ponent of x in the basis A, and 0 otherwise, is called
the basis mapping; it is a linear injective mapping of X
into the space KT of functions on T with values in K.
In this case the image :::( X) consists of all functions
with a finite number of non-zero values (functions of
finite support). This interpretation permits one to
define a generalized basis of a vector space X over a
field K as a bijective linear mapping from it to some
subspace K(T) of the space KT of functions on T with
values in K, where T is some suitably chosen set. How-

ever, unless additional restrictions (e.g. an order) and
additional structures (e.g. a topology) are imposed on
T, and corresponding compatible conditions on K (n
are introduced, the concept of a generalized basis is sel-
dom of use in practice.
A basis of a vector space X is sometimes called an
algebraic basis; in this way it is stressed that there is no
connection with additional structures on X, even if they
are compatible with its vector structure.
A Hamel basis is a basis of the field of real numbers
R, considered as a vector space over the field of
rational numbers. It was introduced by G. Hamel [4] to
obtain a discontinuous solution of the functional equa-
tion f(x +y)= f(x)+ f(Y); the graph of its solution is
everywhere dense in the plane R2. To each almost-
periodic function corresponds some countable Hamel
basis /3 such that each Fourier exponent An of this
function belongs to the linear envelope of /3. The ele-
ments of /3 may be so chosen that they belong to a
sequence {Ai}; the set /3 is said to be a basis of the
almost-periodic functions. An analogous basis has been
constructed in a ring containing a skew-field P and
which has the unit of P as its own unit. An algebraic
basis of an arbitrary vector space is also sometimes
referred to as a Hamel basis.
A topological basis (a basis of a topological vector
space X over a field K) is a set A ={at : tET} eX with
properties and functions analogous to those of the alge-
braic basis of the vector space. The concept of a topo-
logical basis, which is one of the most important ones
in functional analysis, generalizes the concept of an
algebraic basis with regard to the topological structure
of X and makes it possible to obtain, for each element
X, its decomposition with respect to the basis {at},
which is moreover unique, i.e. a representation of x as
a limit (in some sense) of linear combinations of ele-
ments at:
x = lim ~~t(x)at,
where ~t(x) are linear functionals on X with values in
K, called the components of x in the basis A, or the
coefficients of the decomposition of x with respect to the
basis A. Clearly, for the decomposition of an arbitrary
x to exist, A must be a complete set in X, and for such
a decomposition to be unique (i.e. for the zero element
of X to have all components equal to zero), A must be
a topologically free set in X.
The sense and the practical significance of a topolog-
ical basis (which will be simply denoted as a 'basis' in
what follows) is to establish a bijective linear mapping
of X, called the basis mapping, ::: into some (depending
on X) space K (n of functions with values in K,
defined on a (topological) space T, viz.:
E(x):x E X ~ ~At) E K(T),
BASIS
where ~x<t)=~t(x), so that, symbolically,
{~t(X)}=K(n and {~X<T)}=X. Owing to its concrete,
effective definition, the structure of K (n is simpler
and more illustrative than that of the abstractly given
X. For instance, an algebraic basis of an infinite-
dimensional Banach space is not countable, while in a
number of cases, if the concept of a basis is suitably
generalized, the cardinality of T is substantially smaller,
and K (n simplifies at the same time.
The space K(n contains all functions of finite sup-
port, and the set of elements of the basis {at} is the
bijective inverse image of the set of functions {~t(s)}
with only one non-zero value which is equal to one:
at = E-l[~t(s)l,
where ~t(s)=l if t=s, and ~t(s)=O if t=l=s. In other
words, at is the generator of a one-dimensional sub-
space At which is complementary in X to the hyper-
plane defined by the equation Mx)=O.
Thus, the role of the basis {at} is to organize, out of
the set of components ~t(x) which constitute the image
of x under the basis mapping, a summable (in some
sense) set {~t(x)at}, i.e. a basis 'decomposes' a space X
into a (generalized) direct sum of one-dimensional sub-
spaces:
A basis is defined in a similar manner in vector spaces
with a uniform, limit (pseudo-topological), linear (L-),
proximity, or other complementary structure.
Generalizations of the concept of a basis may be and
in fact have been given in various directions. Thus, the
introduction of a topology and a measure on T leads to
the concept of the so-called continuous sum of elements
from X and to corresponding integral representations;
the decomposition of the space X into (not necessarily
one-dimensional) components is used in the spectral
theory of linear operators; the consideration of arbi-
trary topological algebras over a field K (e.g. algebras
of measures on T with values in K or even in X, alge-
bras of projection operators, etc.) instead of K (n
makes it possible to concretize many notions of
abstract duality for topological vector spaces and, in
particular, to employ the well-developed apparatus of
the theory of characters.
A countable basis, which is the most extensively stu-
died and, from the practical point of view, the most
important example of a basis, is a sequence {ai} of ele-
ments of a space X such that each element x is in
unique correspondence with its series expansion with
respect to the basis {ai}
~~i(x)ai' ~i(X)EK,
which (in the topology of X) converges to x. Here,
T=Z, and there exists a natural order in it. A count-
able basis is often simply called a 'basis'. A weak count-
347
BASIS
able basis is defined in an analogous manner if weak
convergence of the expansion is understood. For
instance, the functions e ikt , k EZ, form a basis in the
spaces Lp ' I <p < 00 (periodic functions absolutely
summable of degree p); on the contrary, these functions
do not form a basis in the spaces L I, L 00 (measurable
functions which almost everywhere coincide with
bounded functions) or C I (continuous periodic func-
tions). A necessary, but by far not sufficient, condition
for the existence of a countable basis is the separability
of X (e.g. a countable basis cannot exist in the space of
measurable functions on an interval [a, b] with values
in R). Moreover, the space 100 of bounded sequences,
not being separable in the topology of 100 , has no
countable basis, but the elements ai = {~ik}, where
~ik = I if i =k, and ~ik =0 if i=Fk, form a basis in the
weak topology a(loo, II)' The question of the existence
of a countable basis in separable Banach spaces (the
basis problem) has been negatively solved [6]. The
analogous problem for nuclear spaces also has a nega-
tive solution [7].
A countable basis is, however, not always 'well-
suited' for applications. For example, the components
~t(x) may be discontinuous, the expansion of x need
not converge unconditionally, etc. In this connection
one puts restrictions on the basis or introduces general-
izations of it.
A basis of countable type is one of the generalizations
of the concept of a countable basis in which, although
T is not countable, nevertheless the decomposition of
x EX with respect to it has a natural definition: the
corresponding space K (T) consists of functions with
countable support. For instance, a complete orthonor-
mal set {ar) in a Hilbert space H is a basis; if x EH,
then ~t(x)= <x, at> (where <','> is the scalar pro-
duct in H) for all (except possibly a countable set of)
indices t E T, and the series ~~Ial converges to x. The
basis mapping is determined by the orthog() lal projec-
tions onto the closed subspaces generated oy the ele-
ments al . A basis of the space AP of all complex-valued
almost-periodic functions on R consists of the functions
eilA; here, T=R, K(T) is the set of countably-valued
functions, and the basis mapping is defined by the for-
mula: fT
:=:[x (A)] = r-----ox
lim -2 1 f x (A)e IlA dA.
T T
An unconditional basis is a countable basis in a space
X such that the decomposition of any element x con-
verges unconditionally (i.e. the sum of the series does
not change if an arbitrary number of its terms is rear-
ranged). For instance, in Co (sequences converging to
zero) and lp (sequences sum mabIe of degree p,
I ~p < (0) the elements a, = {8,k} form an uncondi-
tional basis; in the space C[a, b] of continuous func-
tions on the interval [a, b] no basis can be uncondi-
tional. An orthonormal countable basis of a Hilbert
space is an unconditional basis. A Banach space with
an unconditional basis is weakly complete (accordingly,
it has a separable dual space) if and only if it contains
no subspace isomorphic to Co (or, correspondingly, II)'
Two bases {ai} and {bi } of the Banach spaces X and
Y, respectively, are said to be equivalent if there exists a
bijective linear mapping T: ai~bi that can be extended
to an isomorphism between X and Y; these bases are
said to be quasi-eqUivalent if they become equivalent as
a result of a certain rearrangement and normalization
of the elements of one of them. In each of the spaces,
11,1 2 , Co all normalized unconditional bases are
equivalent. However, there exist normalized bases not
equivalent to orthonormal ones.
A summable basis - a generalization of the concept
of an unconditional basis corresponding to a set T of
arbitrary cardinality and becoming identical with it if
T=Z - is a set A ={at : tET} such that for an arbi-
trary element x EX there exists a set of linear combina-
tions (partial sums) of elements from A, which is called
a generalized decomposition of x, which is summable to
x. This means that for any neighbourhood U ex of
zero it is possible to find a finite subset Au eA such
that for any finite set A' =:J A u the relation
[ ~,~tat - x] EO U,
tEA
is true, i.e. when the partial sums form a Cauchy sys-
tem (Cauchy filter). For instance, an arbitrary ortho-
normal basis of a Hilbert space is a summable basis. A
weakly summable basis is defined in a similar way. A
totally summable basis is a summable basis such that
there exists a bounded set B for which the set of semi-
norms {pB(~tal)} is summable. A totally summable
basis is at most countable. In a dual nuclear space all
weakly summable bases are totally summable.
An absolute basis (absolutely summable basis) is a
summable basis of a locally convex space over a
normed field such that for any neighbourhood U of
zero and for each t E T the family of semi-norms
{p u( al )} is summable. All unconditional countable
bases are absolute, i.e. the series ~ I ~Jx) Ip(ai)
converges for all x EX and all continuous semi-norms
pC} Of all Banach spaces only the space II has an
absolute countable basis. If a Frechet space has an
absolute basis, all its unconditional bases are absolute.
In nuclear Frechet spaces any countable basis (if it
exists) is absolute [13].
A Schauder basis is a basis {a l : t E T} of a space X
such that the basis mapping defined by it is continuous
(and is therefore an isomorphism onto some space
 BATEMAN METHOD

K(T), i.e. a basis in which the components ~t(x) for (1974),304-313. (Funktsional. Analiz. i Prilozhen. 8, no. 4
any x EX and, in particular, the coefficients of the (1974), 35-47)
[8] EDWARDS, R.E.: Functional analysis: theory and applications.,
decomposition of x with respect to this basis, are con- Holt, Rinehardt, Winston, 1965.
tinuous functionals on X. This basis was first defined [9] DIEUDONNE, J.: 'Sur les espaces de Kothe', J. d'Analyse Math.
by 1. Schauder [5] for the case T=Z. The concept of a 1 (1951), 81-115.
[10] ARSOVE, M.G.: 'The Paley-Wiener theorem in metric linear
Schauder basis is the most important of all modifica-
spaces', Pacific J. Math. 10 (1960), 365-379.
tions of the concept of a basis. [II] BESSAGA, C. and PELCZYNSKl, A.: 'Spaces of continuous func-
A Schauder basis is characterized by the fact that tions IV', Studia Math. 19 (1960),53-62.
{at} and {~t} form a bi-orthogonal system. Thus, the [12] JAMES, R.C.: 'Bases and reflexivity in Banach spaces', Ann. of
Math. (2) 52, no. 3 (1950), 518-527.
sequences ai = {t3 ik } form countable Schauder bases in [13] DYNIN, A. and MITYAGIN, B.: 'Criterion for nuclearity in
the spaces Co and ~, p ~ 1. A countable Schauder basis terms of approximate dimension', Bull. A cad. Polon. Sci. ser.
forms a Haar system in the space C[a, b]. In complete Sci. Math., Astr. Phys. 8 (1960), 535-540.
[14] DAY, M.M.: Normed linear spaces, Springer, 1958.
metric vector spaces (in particular, in Banach spaces)
[15] PIETSCH, A.: Nuclear locally convex spaces, Springer, 1972
all countable bases are Schauder bases [10]. In Frechet (translated from the German).
spaces the concept of a weak basis and a Schauder [16] SINGER, I.: Bases in Banach spaces, 1-2, Springer, 1970-1981.
basis are identical [11]. In barrelled spaces in which M.l. Vortsekhovskir
there are no linear continuous functionals, there is also M.l. Kadets
no Schauder basis [8]. However, if a weak Schauder AMS 1980 Subject Classification: 08-XX, 10l05,
basis exists in these spaces, it is an ordinary Schauder 46815, 46A35
basis [9]. A barrelled locally convex space with a count-
able Schauder basis is reflexive if and only if this basis BATEMAN FUNCTION, kjunction - The function
is at the same time a shrinking set, i.e. if the {~t} 1T/2

corresponding to it will be a basis in the dual space X kv(x) = 1. f cos(x tg8-p8)d8, (I)
'IT 0
and will be boundedly complete, i.e. if the boundedness
where x and /l are real numbers. The function was
of the set of partial sums of a series ~~iai implies that
I defined by H. Bateman [1]. The Bateman function may
this series is convergent [12]. If a Schauder basis is an be expressed in the form of a confluent hypergeometric
unconditional basis in a Banach space, then it is a function of the second kind \{tea, h, x):
shrinking set (or a boundedly complete set) if and only
if X does not contain subspaces isomorphic to II (or,
[(p+l)k 2v (x) = e- x 'I'(-p,O;2x). x>O. (2)

respectively, to co). The relation (2) is conveniently taken as the definition

A Schauder basis in a locally convex space is
equicontinuous if for any neighbourhood U of zero it is
possible to find a neighbourhood V of zero such that
I ~,(x) Ipu(a,) ~ pv(x)
for all x EX, t E T. All Schauder bases of a barrelled
space are equicontinuous, and each complete locally
convex space with a countable equicontinuous basis can
be identified with some sequence space [15]. An
equicontinuous basis of a nuclear space is absolute.
References
[1] COHN, P.M.: Universal algebra, Reidel, 1981.
[2] MAL'TSEV, A.I.: Algebraic systems, Springer, 1973 (translated
from the Russian).
[3] BOURBAKI, N.: Elements of mathematics. Algebra: Algebraic
structures. Linear algebra, 1, Addison-Wesley, 1974,
Chapt. 1; 2 (translated from the French).
[4] HAMEL, G.: 'Eine Basis aller Zahlen und die unstetigen
Losungen der Funktionalgleichung: f(x +y) =f(x) +f(Y)',
Math. Ann. 60 (1905), 459-462.
[5] SCHAUDER, J.: 'Zur Theorie stetiger Abbildungen in
Funktionalraumen', Math. Z. 26 (1927), 47-65; 417-431.
[6] ENFLO, P.: 'A counterexample to the approximation problem
in Banach spaces', Acta Math. 130 (1973),309-317.
[7] ZOBIN, N.M. and MITYAGIN, B.S.: 'Examples of nuclear linear
metric spaces without a basis', Functional Anal. Appl. 8, no. 4
of the Bateman function in the complex plane with the
cut (- 00, 0]. The following relations are valid: for case
(1)
kvC - x) = Lv(x),
for case (2)
k2v(-~iO) = L2vW-e='=""'e~CP(-p.O; -2~).
where ~>O, and !lJ(a, b ; x) is a confluent hyper-
geometric function of the first kind.
References
[I] BATEMAN. H.: 'The k -function, a particular case of the con-
fluent hypergeometric function', Trans. A mer. Math. Soc. 33
(1931),817-831.
[2] BATEMAN, H. and ERDELYI, A. (EDs.): Higher transcendental
functions, I, McGraw-Hill. 1953.
L.N. Karmazina
AMS 1980 Subject Classification: 33A30
BATEMAN METHOD - A method for approximating
the integral operator of a one-dimensional integral
Fredholm equation of the second kind; it is a particular
case of the method of degenerate kernels (cf. Degen-
erate kernels, method of). In Bateman's method, the
degenerate kernel KN(X, s) is constructed according to
the rule:

349
BATEMAN METHOD

KN(x, S) = Editorial comments.

o K(x, Sl)
References
K(xJ, S) K(xJ, Sl) [A1] LIPTSER, R.S. and SHIRYAYEV, A.N.: Statistics of random
processes, 1, Springer, 1977, Section 7.9.

AMS 1980 Subject Classification: 62A 15

BAYESIAN APPROACH to statistical problems - An

approach based on the assumption that to any parame-
ter in a statistical problem there can be assigned a
definite probability distribution. Any general statistical
where Si, Xi, i = I, ... ,N, are certain points on the decision problem is determined by the following ele-
integration segment of the integral equation considered. ments: by a space (X, f!4 x) of (potential) samples X, by
The method was proposed by H. Bateman [I]. a space (8, f!4 e) of values of the unknown parameter
References B, by a family of probability distributions {Po: BEe}
[I] BATEMAN, H.: Messeng. Math. 37 (1908), 179-187.
on (X, !J8 x), by a space of decisions (D, f!4 D) and by a
[2] KANTORovICH, L.V. and KRYLOV, V.I.: Approximate methods
of higher analysis, Noordhoff, 1958 (translated from the Rus- function L(B, d), which characterizes the losses caused
by accepting the decision d when the true value of the
sian). A.B. Bakushinskit
parametar is B. The objective of decision making is to
Editorial comments. find in a certain sense an optimal rule (decision func-
References tion) o=o(x), assigning to each result of an observation
[A1] BATEMAN, H.: Proc. Roy. Soc. A (1922),441-449.
XEX the decision o(x)ED. In the Bayesian approach,
AMS 1980 Subject Classification: 45805, 45L05 when it is assumed that the unknown parameter B is a
random variable with a given (a priori) distribution
BAYES FORMULA - A formula with which it is pos-
7T=7T(d8) on (8, 88 e) the best decision function
sible to compute a posteriori probabilities of events (or (Bayesian decision function) o = o (x) is defined as the
of hypotheses) from a priori probabilities. Let function for which the minimum expected loss
AI, ... ,An be a complete group of incompatible
inf,s p(7T, 0), where
events: UAi=n, Ai nAJ== 0 if i=/=j. Then the a
posteriori probability P(A i I B) of event A, if given that p(7T, 8) = j p(O, 8)7T(dO),
EI
event B with P(BO has already occurred may be
and
found by Bayes' formula: p(O, 8) = jL(0.8(xP e(dx)
.\
P(A,)P(B I A,)
peA, I B) = n
(*) is attained. Thus,
~ peA, )P(B I A,)
p(7T.8*) = inf j j L(O, 8(x Pe(dX)7T(dO).
ii (-J X
where P(Ai) is the a priori probability of Ai, PCB IAi) is
the conditional probability of event B occurring given In searching for the Bayesian decision function
event A, (with P(AiO) has taken place. The formula 0* = 0* (x), the following remark IS useful. Let
was demonstrated by T. Bayes in 1763. ) Po(dx) =p(x I B) d,u(x), w(dB) = 7T( B) dv( B). where ,u and v
Formula (*) is a special case of the following abstract are certain a-finite measures. One then finds, assuming
variant of Bayes' formula. Let B and ~ be random ele- that the order of integration may be changed.
ments with values in measurable spaces (8, BE) and j jL(0,8(xP o(dx)7T(dO) =
(X,B x ) and let EIK(B)I<oo. Put for any set (oJ .\

A EF~=a{w: ~(w)}. = jjL(0.8(xp(xI0)7T(0)dp.(x)dv(0) =

(oJ .\
G(A) = jg(0(wE[14(w) I Fel(w)P(dw).
11

where F8=a{w: B(w)} and JA(w) is the indicator of the

= f.\ dp.(x) lr(j-) L(O, 8(x)p(x I O)7T(O)dV(O)l

set A. Then the measure G is absolutely continuous It is seen from the above that for a given x EX, 0* (x) is
with respect to the measure P (CP) and that value of d* for which
E[g(B) I Fd(w)=(dG / dP)(w). where (dC / dP)(w) is the
Radon - Nikodym derivative of C with respect to P. in: j L(O. d)p(x 10)7T(O)dv(O)
d (_)

References is attained, or. what is equivalent for which

[I J KOl.MOGORm. 1\.1\.: FOIlI1c/mio/ls oj (he (heorl' "/pl'lllwlllllll.
Chelsea. reprint. 1950 (translated fmIll the Ru"ian). inf f L(O, d) p(xp(x)
10)77(0) dv(O).
A. S. Sh /'"),(1('1' d H

3S0

is attained, where
p(X) ::: [p(X I O)TT(O) dp(O).
But, according to the Bayes fonnula
(L(O, d) p(x IO)TT(O) dv(O) ::: E[L(O, d) Ix].
~ p(x)
Thus, for a given x, 13* (x) is that value of d* for which
the conditional average loss E[L(O, d) Ix] attains a
minimum.
Example. (The Bayesian approach applied to the case
of distinguishing between two simple hypotheses.) Let
8={0],02}, D={d l ,d2 }, Lij=L=(Oi,dj ), i,j=I,2;
7T(Od=7TJ, 7T(02) = 7T2, 7TI +7T2 = 1. If the solution di is
identified with the acceptance of the hypothesis Hi:
O=Oi' it is natural to assume that LII <L 12 , L 22 <L 21
Then
p(TT,O) ::: j[TTIP(X IOI)L(OI, o(x+
x unconditional distribution of Y.
+TT2P(X I(2)L(02, o(x]d}L(x)
implies that infs p( 7T, (3) is attained for the function
o'(X) :::
The advantage of the Bayesian approach consists in
the fact that, unlike the losses p(O, (3), the expected
losses p(7T, (3) are numbers which are dependent on the
unknown parameter 0, and, consequently, it is known
that solutions 13; for which
p(TT, 0;) .;;; inf p(TT, 0)+(,
6
and which are, if not optimal, at least (-optimal (>0),
are certain to exist. The disadvantage of the Bayesian
approach is the necessity of postulating both the
existence of an a priori distribution of the unknown
parameter and its precise form (the latter disadvantage
may be overcome to a certain extent by adopting an
empirical Bayesian approach, cf. Bayesian approach,
empirical).
References
[1] WALD, A.: Statistical decision functions, Wiley, 1950.
[2] GROOT, M.H. DE: Optimal statistical decisions, McGraw-Hill,
1970.
A.N. Shiryaev
AMS 1980 Subject Classification: 62A15, 62C10
BAYESIAN APPROACH, EMPIRICAL - A statistical
interpretation of the Bayesian approach yielding conclu-
sions on unobservable parameters even if their a priori
distribution is unknown. Let (Y, X) be a random vector
for which the density p (y Ix) of the conditional distri-
bution of Y for any given value of the random parame-
BA YESIAN APPROACH, EMPIRICAL
ter X=x is known. If, as a result of some experiment,
only the realization of Y is observed, while the
corresponding realization of X is unknown, and if it is
necessary to estimate the value of a given function cJ>(X)
of the non-observed realization, then, in accordance
with the empirical Bayesian approach, the conditional
mathematical expectation E{ cJ>(X) I Y} should be used
as an approximate value tf(Y) for cJ>(X). In view of the
Bayes fonnula, this expectation is given by the formula
Y - jcp(x)p(Ylx)p(x)d}L(x) (1)
1/;( ) - q(Y) ,
where
q(y)::: jP(Ylx)p(x)d}L(x), (2)
p (x) is the density of the unconditional (a priori) distri-
bution of X, }L(x) is the corresponding a-finite measure;
and the function q(y) represents the density of the
If the a priori density p (x) is unknown, it is not pos-
sible to compute the values of tf and q. However, if a
sufficiently large number of realizations of the random
variables Y I, . . . , Yk , which are drawn from the distri-
bution with density q(y), is known, it is possible to
construct a consistent estimator q(y), which depends
only on Y I , . . . , Y k S.N. Bernshteln [I] proposed to
estimate the value of tf( Y) by substituting q(y) for q (y)
in (2), finding the solution p(x) of this integral equa-
tion, and then substituting p
and q
in the right-hand
side of (1). However, this method is difficult, since solv-
ing this integral equation (2) is an ill-posed problem in
numerical mathematics.
In certain special cases the statistical approach may
be employed not only to estimate q, but also tf [3]. This
is possible if the identity
cp(x)p(y Ix) ::: A(y)r[z(y) Ix], (3)
involving x and y, is true. In (3), '\(y) and z (y) are
functions which depend on y only, while r(z Ix), being
a function of z, is a probability density (i.e. may be
regarded as the density of an arbitrary distribution of
some random variable Z for a given value X=x). If (3)
is true, the numerator of (I) is equal to the product
'\(Y)s [z (Y)], where s(z)= fr(z Ix)p(x) d}L(x) is the
density of the unconditional distribution of Z. Thus, if
a sufficiently large number of realizations of indepen-
dent random variables Z I, . . . ,Zm with density distri-
bution s (z) is available, then it is possible to construct
a consistent estimator s(z) Jor
s(z), and hence also to
find a consistent estimator tf( Y) for tf( Y):
CP(X) ~ I/;(Y) ~ ~Y) ::: A(Y!S[Z(Y)]. (4)
q(Y)
For instance, if one has to estimate tf(X)=xh, where
h is a positive integer, and p(y Ix)=xYe- x Iy!
351
BA YESIAN APPROACH, EMPIRICAL
(y=o, I, ... ; x>O), then cp(x)p(y IX)=A(y)p(y+h Ix),
where A(y)=(y +h)! IY!. Since, here, r(z Ix)=p(z Ix),
one has s (z) = q (z). Accordingly,
~(y) = A( Y)q( Y + h) I q( y), i.e. only the sequence of
realizations Y I, Y 2, . . . , is required to find 1/;. If, on
the other hand, p(Ylx)=b(Ylx)=C{xY(l-x)n- y
(y=O, ... ,n; n is a positive integer; O':;;;x.:;;;l,
C{ =<;, then 1/;(x)p(y I x)=A(y)r(y +h I x), where
A(Y)= C{ I C{ ti and r(z Ix)=b n +h(z Ix)*p(z Ix). For
this reason two sequences of empirical values Y i and Z}
are required in this case to construct 1/;( Y).
This form of the empirical Bayesian approach is
applicable to the very narrow class of densities p (y Ix)
and functions cp(x) which satisfy condition (3); even if
this condition is in fact met, the construction of the
estimator (4) is subject to the observability of the ran-
dom variables Zj' the distribution of which usually
differs from that of the variables Yi which are observed
directly. For practical purposes, it is preferable to use
the empirical Bayesian approach in a modified form, in
which these disadvantages are absent. In this modifica-
tion the approximation which is constructed does not
yield a consistent estimator of 1/;(y) (such an estimator
may even be non-existent), but rather upper and lower
estimators of this function, which are found by solving
a problem of linear programming, as follows. Let
'II I (Y) and 'II2(Y) be the constrained minimum and
maximum of the linear functional (with respect to the
unknown a priori densi ty p (x in the n umera tor of (1),
calculated under the linear constraints p(x);;?O,
jp (x) dJL(x) = I and q(Y) jp(Ylx)p(x)dJL(x)=q(y),
where q( Y) is the estimator of q (Y) mentioned above,
constructed from the results of the observations
Y I, . . . , Yk . One may conclude in such a case that
'III (Y) I q(Y).:;;;1/;(Y).:;;;'lI 2 (Y) I q(Y), where the proba-
bility of the truth of this conclusion tends to one (by
virtue of the law of large numbers) as the number of
random variables Yi , used to construct thl) estimator
q(Y), increases without limit. Other modifications of
the empirical Bayesian approach are also possible -
for example, by adding to the last-named condition
q(Y)=q(Y) a finite number of conditions of the form
q (Yi) = q(Yi), where y/ are preliminarily given numbers;
q
if is replaced by the corresponding confidence bounds
for q, the conditions are obtained in the form of ine-
qualities ql(YJ':;;;q(Yi)':;;;Q2(Y/). etc.
In certain cases. which are important in practice.
satisfactory majorants. which can be computed without
the use of the laborious method of linear programming.
can be found for the functions 'II I and 'lI2 (see the
example in the entry Sample method whieh deals with
statistical control).
See the entry Discriminant analysis for the applica-
352
tions of the empirical Bayesian approach to hypotheses
testing concerning the values of random parameters.
References
[1] BERNSHTEiN, S.N.: 'On 'fiducial' probabilities of Fisher', lzv.
Akad. Nauk SSSR Ser. Mat. 5 (1941),85-94 (in Russian).
English abstract.
[2] BOL'sHEv, L.N.: 'Applications of the empirical Bayes
approach', in Internat. Congress Mathematicians Nice, 1970,
Vol. 3, Gauthier-Villars, 1971, pp. 241-247.
[3] ROBBINS, H.: 'An empirical Bayes approach to statistics', in
Proc. Berkeley Symp. Math. Statist. Probab., Vol. 1, Berkeley-
Los Angeles, 1956, pp. 157-163.
L.N. Bol'shev
AMS 1980 Subject Classification: 62A 15, 62C15
BAYESIAN DECISION FUNCflON - A rule (func-
tion) ~=~(x) which associates with each result x of a
statistical experiment a decision ~(x) with values in a
given set of decisions, and which attains the minimum
expected loss as defined in the framework of the Bayes-
ian approach to statistical problems.
A.N. Shiryaev
AMS 1980 Subject Classification: 62C10
BAYESIAN ESTIMATOR - An estimator of an
unknown parameter from the results of observations
using the Bayesian approach. In such an approach to
the problem of statistical estimation it is usually
assumed that the unknown parameter 0 E e C Rk is a
random variable with given a priori distribution
7T=7T(dO), that the space of decisions D is identical to
the set e and that the loss L(O, d) expresses the devia-
tion between the variable 0 and its estimator d. It is
therefore supposed, as a rule, that the function L (0, d)
has the form L(O, d)=a(O)A(O-d), where A is some
non-negative function of the error vector 8-d. If k = 1,
it is often assumed that A(O-d)= 18-d I", 0:>0; the
most useful and mathematically the most convenient is
the quadratic loss function L(O, d)= 18-d 12. For such
a loss function the Bayesian estimator (Bayesian deci-
sion function) 0* = 0* (x) is defined as the function for
which the minimum total loss
inf p(?T. 8) = inf f f 10-8(x) 12 Po (dx)?T(dO),
8 8 (_) X
IS attained. or, equivalently. for which the minimum
conditional loss
infE{[0-8(x)f I x)
Ii
is attained. It follows that in the case of a quadratic
loss function the Bayesian estimator 0* (x) coincides
with the a posteriori average 0*(x)=E(8Ix). and the
Bayes risk is
p(?T,8') = E[O(Oix)j,
where 0(01 xl is the variance of the a posteriori distri-
bution:
O(Olx) = E{[O-E(Olx)flx}.
Example. Let x = (x 1 . . , x n ), where XI, . . . 'X n are

independent identically-distributed random variables
with normal distributions N(), a2), a2 is known, while
the unknown parameter () has the normal distribution
N (JJ., ~). Since the a posteriori distribution for ()
(where x is given) is normal N(JJ.n, ,.~) with
- -2 +
_ nxo
JLn - -2
fJ.T
-2
-2'
-2 _
Tn - no
-2 + T-2,
no +,.
where x = (x I + ... + Xn) / n, it follows that for the
quadratic loss function 1 ()-d 12 the Bayesian estima-
tor is 5' (X) =JJ.n, while the Bayesian risk IS
,.~ =0'2~ / (n~ +0'2).
A.N. Shiryaev
Editorial comments.
References
[A1] SVERDRUP, E.: Laws and chance variations, 1, North-Holland,
1967, Chap!. 6, Section 4.
AMS 1980 Subject Classification: 62F15
BEHNKE - STEIN THEOREM - The umon of
domains of holomorphy Gk C en , k = I, 2, ... , for
which Gk C Gk + I for all k, is also a domain of holo-
morphy. The Behnke - Stein theorem is valid not only
in the complex Euclidean space en, but also on any
Stein manifold. If the sequence Gk is not monotone
increasing by imbedding, the theorem is not valid; e.g.
the union of the two domains of holomorphy
G1 = {(ZI,Z2): I ZI I <I, I Z2 I <2}
and
G2 = {(zJ, Z2): I ZI I <2, I Z2 I <I}
in C2 is not a domain of holomorphy.
References
[1] BEHNKE, H. and STEIN, K.: 'Konvergente Folgen von
Regularitatsbereichen und die Meromorphiekonvexitlit', Math.
Ann. 116 (1938),204-216.
[2] VLADlMIROV, V.S.: Methods of the theory offunctions of several
complex variables, M.LT., 1966 (translated from the Russian).
E.M. Chirka
AMS 1980 Subject Classification: 32005
BEHRENS- FISHER PROBLEM - An analytical
problem which arose in the context of the statistical
problem of comparing, starting from empirical data, the
mathematical expectations of two normal distributions,
the variances of which are unknown (it is assumed that
the ratio of the variances is also unknown). This prob-
lem was posed by W.U. Behrens [l] in connection with
processing crop data. The modern formulation of the
Behrens - Fisher problem is due to R. Fisher and is
based on the concept of sufficient statistics. Let
XII, ... ,X ln , and X 21 , ... ,X2n , be mutually
independent random variables with a normal distribu-
tion, and let EX li =JJ.J, E(Xli - JJ.d =O'f (i = I, ... ,n I)
and EX2j=JJ.2' E(X2j-JJ.d=0'~ (j=I, ... ,n2). It is
assumed that the values of the mathematical expect a-
BEHRENS - FISHER PROBLEM
tions JJ.l, JJ.2, of the variances aT, o'~ and of their ratio
aT / o'~ are unknown. A sufficient statistic in the case
nl, n2~2 is a four-dimensional vector (~\, X2, ST, sh
the components of which are expressed by the formulas
n, n2_
Sy = ~(Xli-:\\)2, S~ = ~(X2)-X2)2,
i=1 )=1
and which are mutually indeJ>endent random variables;
~(XI-JJ.I)/O'l and yr;;;(X2 -JJ.2)/(J2 have a
standard normal distribution, while Sf / (JT and
S~ / (J~ have a 'chi-squared' distribution with n I - I
and n2 -} degrees of freedom, respectively. Since a suf-
ficient statistic contains the same information about the
unknown parameters JJ.I, JJ.2, (JT, (J~ as the initial
n I + n 2 random variables X Ii and X 2j' it follows that
only the sufficient statistics need be considered in test-
ing hypotheses about the values of these parameters.
In particular, this idea is the basis of the modern for-
mulation of the problem of hypotheses testing, concern-
ing the hypothesis JJ.I - JJ.2 =~, where ~ is a previously
given number; here the Behrens- Fisher problem is
reduced to finding a set Ka in the space of possible
values of the random variables XI - X 2, ST, S~ such
that, if the hypothesis being tested is correct, the proba-
bility of the event (XI-X2,Sf,S~)EKa does not
depend on all the unknown parameters and is exactly
equal to a given number a in the interval O<a< l.
The question of the existence of a solution to the
Behrens - Fisher problem was discussed at length by
prominent mathematicians (mainly in connection with
the approach to the problem taken by R.A. Fisher,
which passed beyond the borders of probability
theory). It was shown by Yu.V. Linnik et aI., in }964,
that if the sample sizes n I and n 2 are of different pari-
ties, a solution Ka to the Behrens- Fisher problem
exists. If the parities of n I and n2 are equal, the
existence of a solution remains an open question.
The Behrens- Fisher problem has often been gen-
eralized and modified. A. Wald, in particular, posed the
problem of finding a set K a in the sample space of the
two variables (XI - X 2) / ST and ST / S~. The question
of the existence of a solution to this problem remains
open. However, it is effectively possible to construct a
set K*,. such that if the hypothesis JJ.I - JJ.2 = ~ being
tested is in fact correct, the probability of the event
X 1 -X 2)/ST,ST /S~)EK:, while still depending on
the unknown ratio aT / O'~, will deviate from the given
a only by a small amount. This fact is the basis of
modern recommendations for the practical construction
of tests to compare JJ.l and JJ.2. Simple tests for the
comparison of JJ.l with JJ.2, which are also computation-
ally convenient, were proposed by V.T. RomanovskiI,
353
BEHRENS - FISHER PROBLEM
M. Bartlett, H. Scheffe and others. However, the statis-
tics of these tests are not expressed in terms of suffi-
cient statistics and are, for this reason, usually less
powerful than tests based on the solution of the
Behrens- Fisher problem and its generalizations.
References
[I] BEHRENS, W.U.: Landwirtsch.lahresber. 68, no. 6 (1929),807-
837.
[2] LiNNIK, YU.V.: Statistical problems with nuisance parameters,
Amer. Math. Soc., 1968 (translated from the Russian).
[3] LINNIK, Yu.V., ROMANOYSKIi, LV. and SUDAKOY, V.N.: 'A
nonrandomized homogeneous test in the Behrens - Fisher
problem'. Soviet Math. Dok!. 5, no. 2 (1964),570-572. (Dok!.
Akad. Nauk SSSR 155, no. 6 (1964), 1262-1264)
L.N. Bol'shev
Editorial comments.
References
[A 1] LiNNIK, YU.V.: 'Randomized homogeneous tests for the
Behrens - Fisher problem', Selected Transl. in Math. Stat.
and Probab. 6 (1966), 207-217 (translated from the Rus-
sian).
AMS 1980 Subject Classification: 62FXX
BELL-SHAPED GAME - A game on the unit square
whose pay-off function takes the form cf>(x - y), where cp
is a positive analytic proper P6lya frequency function,
1.e. :
1) cp(u) is defined for all UE(-oo, 00);
2) for any n and any sets - 00 <x 1< ... <xn < 00
and - 00 <y I < ... <Yn < 00 there is an inequality
det II cf>(x, - Yj) II ~O;
3) for any set {xd (correspondingly, {Yd) there is a
set {Yd (correspondingly, {Xk }) such that
det II cf>(Xj - Yj) II >0;
4) fjJoo
cf>(u) du < 00.
An example of a bell-shaped game is a game with
pay-off function e -(x _y)2. The optimal strategies of
players in a bell-shaped game are unique and are
piecewise-constant distributions with a finite number of
steps. The value of a game with pay-o), function
cf>(?-.(x -y)), as ?-.~OO, moves towards zero, while the
number of points in the supports of the optimal stra-
tegies grows unboundedly.
References
[I] KARl.IN, S.: Mathematical method.1 and theory in Rames, pro-
gramming and economics, Addison-Wesley, 1959.
v.K. Domanskir
AMS 1980 Subject Classification: 90005
BELLMAN EQUATION - I) A partial differential
equation of a special type to solve a problem of optimal
control. If the solution of Cauchy's problem for the
Bellman equation can be found, the optimal solution of
the original problem is readily obtained.
2) A recurrence relation for the solution of a discrete
problem of optimal control. The method for obtaining
354
the optimal solution with the aid of Bellman's equation
is known as dynamic programming.
References
[1] BELLMAN, R.: Dynamic programming, Princeton Univ. Press,
1957.
V. G. Karmanov
Editorial comments. The Bellman equation for
continuous-time optimal control problems is also often
called the dynamic programming equation. Cf. the article
Optimality, sufficient conditions for for examples and more
details. There is also a variant for stochastic optimal control
problems.
References
[A1] FLEMING, W.H. and RISHEL, R.W.: Deterministic and sto-
chastic optimal control, Springer, 1975.
AMS 1980 Subject Classification: 35QXX, 90C39
BELLMAN - IlARR.J.s PROCESS - A special case of
an age-dependent branching process (cf. Branching pro-
cess, age~dependent). It was first studied by R. Bellman
and T.E. Harris [1]. In the Bellman - Harris process it
is assumed that particles live, independently of each
other, for random periods of time, and produce a ran-
dom number of new particles at the end of their life
time. If G (t) is the distribution function of the life
times of the individual particles, if h (s) is the generat-
ing function of the number of direct descendants of one
particle, and if at time t =0 the age of the particle was
zero, then the generating function F(t;s)=Esl'(t) of the
number of particles p,(t) satisfies the non-linear integral
equation
F(t;s) = jh(F(t-u;sdG(u)+s(l-G(t.
o
If
G(t) = l-e- A1 , t?O,
the Bellman - Harris process is a Markov branching
process with continuous time.
References
[I] BELLMAN. R. and HARRIS. T.E.: 'On the theory of age-
dependent stochastic branching processes', Proc. Nat. Acad.
Sci. USA 34 (1948), 601-604.
B.A. Sevast ')'anov
AMS 1980 Subject Classification: 60J80
BELTRAMI COORDINATES - Coordinates on a sur-
face in which all geodesic lines are expressed by linear
equations. Surfaces on which the introduction of Bel-
trami coordinates is possible are said to he projective
(projective transformations of such surfaces convert
geodesics into geodesics). So called after E. Beltrami.
who studied such surfaces in 1866.
E. V S'hikin
AMS 1980 Subject Classification: 53A05
BELTRAMI - ENNEPER THEOREM - A theorem
estahlishing the following property of asymptotic lines

(d. Asymptotic line) on a surface of negative curvature
(E. Beltrami, 1866; A. Enneper, 1870): If the curvature
of an asymptotic line at a given point is non-zero, then
the square of the torsion of this line is equal to the
absolute value of the curvature of the surface at this
point. The theorem can also be applied to the case
when the curvature of the asymptotic line at a given
point is zero. The square of the torsion is replaced by
the square of the rate of rotation of the tangent plane
to the surface at this point during the motion along the
asymptote. Asymptotic lines emanating from the same
point have torsions of equal absolute values and of
opposite signs.
E. V Shikin
Editorial comments. A useful general reference is [A 1].
References
[A1] SPIVAK. A.M.: Differential geometry, 3, Publish or Perish,
1975.
AMS 1980 Subject Classification: 53A05
BELTRAMI EQUATION - See Laplace-Beltrami
equation.
AMS 1980 Subject Classification: 58G05, 58G25
BELTRAMI INTERPRETATION - A realization of a
part of the Lobachevskii plane on a pseudo-sphere - a
surface of constant negative curvature. In the Beltrami
in terpretation, the geodesic lines (cf. Geodesic line) and
their segments on the pseudo-sphere play the role of
straight lines and their segments on the Lobachevskil
plane. The isometric mappings of the pseudo-sphere
onto itself represent the movements in the Lobachevskii
plane. with preservation of horocircles. The lengths,
angles and areas on the pseudo-sphere correspond to
the lengths, angles and areas on the Lobachevskii
plane. Under these conditions, to each theorem of
Lobachevskii planimetry which refers to a part of the
Lobachevskii plane there corresponds an analogous
theorem of the intrinsic geometry of the pseudo-sphere.
The Beltrami interpretation realizes a part of the
Lobachevskii plane, but the entire Lobachevskii plane
cannot be realized in three-dimensional Euclidean
space as a regular surface (Hilbert's theorem).
The Beltrami interpretation was proposed by E. Bel-
trami in 1868 [1]. This publication represented the first
realization of Lobachevskii's 'imaginary geometry' in
three-dimensional Euclidean space.
For other interpretations of the Lobachevskii
geometry see Klein interpretation; Poincare model.
References
[1] BELTRAMI. E.: 'Saggio di interpretazione della geometria non-
euc1idea'. in Opere Mat .. Vol. 1. 1868. pp. 374-405.
[2] KAGAN. V.F.: The foundations of geometry. 1, Moscow-
Leningrad. 1949 (in Russian).
E. V Shikin
BENDIXSON CRITERION
Editorial comments.
References
[A 1] KLEIN, F.: Vorlesungen uber nicht-Euklidische Geometrie,
Springer, 1968.
AMS 1980 Subject Classification: 51 M1 0, 53A35
BELTRAMI METHOD - A method for solving wave
equations with three spatial variables, proposed by E.
Beltrami in 1864. The method is based on the fact that
the wave equation on the surface of the characteristic
cone can be converted with the aid of interior differen-
tial operators to an especially simple form which can be
used for the solution u(x, t)
4'lT t5u(xo, to) = /} [u+to ~~ J do,
where g is the sphere {x: 1 x-xo 1 =to}, da is its sur-
face element and a/ an
denotes differentiation in the
direction of the outward normal on the cone
1 x - x 0 12 = (t - to)2. The method is applicable to the
cases of inhomogeneous equations and to equations
with any odd number of spatial variables.
References
[I] COURANT, R. and HiLBERT, D.: Methods ofmathematicalphy-
sics. Partial differential equations, 2, Interscience, 1965
(translated from the German).
Sh.A. Alimov
VA. II'in
Editorial comments. The formula above is called
Beltrami's formula.
References
[A 1] JOHN, F.: Partial differential equations, Springer, 1978.
AMS 1980 Subject Classification: 35L05, 35C15
BENDIXSON CRITERION - A theorem that permits
one to establish the absence of closed trajectories of
dynamical systems in the plane, defined by the equa-
tion
x' = P(x,y), y' = Q(x,y). (*)
The criterion was first formulated by I. Bendixson [1]
as follows: If in a simply-connected domain G the
expression P~ + Q~ has constant sign (i.e. the sign
remains unchanged and the expression vanishes only at
isolated points or on a curve), then the system (*) has
no closed trajectories in the domain G. This criterion
was generalized by H. Dulac [2] as follows: If G is a
simply-connected domain in the (x, y)-plane, if the
functions P and Q E C 1(G), and if a function
j(X,Y)EC 1(G) can be found such that
f ({Mf!l
ax + 2!J!2l}
ay dx dy =1= 0
for any simply-connected subdomain D C G, then the
domain G does not contain any simple rectifiable
closed curve consisting of trajectories and singular
355
BENDIXSON CRITERION
points of the system (*). If the domain G is an annulus,
a similar theorem states that a closed trajectory of (*),
if it exists, is unique. A generalization applying to the
case of system (*) with cylindrical phase space [3] is
also possible.
References
(I] BENDIXSON, I.: 'Sur les courbes definies par des equations
differentielles'. Acta Math. 24 (1901), 1-88.
[2] DULAC, H.: 'Recherches des cycles Iimites. C.R. Acad. Sci.
Paris Ser. 1 Math. 204 (1937). 1703-1706.
[3] ANDRONOY, A.A., VITI, AA. and KHAiKIN, S.E.: Theory of
oscillations, Moscow, 1959 (in Russian).
N.Kh. Rozov
Editorial comments. Bendixson's criterion is also called
the Poincare- Bendixson theorem.
References
[A 1] ANDRONOY, A.A, LEONTOYICH, E.A., GORDON, 1.1. and
MAIER, A.G.: Qualitative theory of second-order dynamic
systems, Wiley, 1973 (translated from the Russian).
AMS 1980 Subject Classification: 34C05, 34C25
BENDIXSON SPHERE - The sphere in real analysis
which is known as the Riemann sphere in the theory of
functions of a complex variable.
Let ~: X2 + y2 + Z2 = I be the unit sphere in the
Euclidean (X, Y, Z)-space, and let N(O, 0, I) and
S (0, 0, - I) be its north and south pole, respectively;
let p and a be planes tangent to ~ at the points Nand
S respectively; let xSy and uNv be coordinate systems
in a and p with axes parallel to the corresponding axes
of the system XOY in the plane Z=O and pointing in
the same directions; let n be the stereographic projec-
tion of ~ onto a from the centre N, and let II' be the
stereographic projection of ~ onto p from the centre S.
Then ~ is the Bendixson sphere with respect to anyone
of the planes a, P. It generates the bijection <p=n'n- 1
of the plane a (punctured at the point S) onto the plane
P, which is punctured at the point N. This bijection is
employed in the study of the behaviour of the trajec-
tories of an autonomous system of real algf) Iraic ordi-
nary differential equations of the second order in a
neighbourhood of infinity in the phase plane (the
right-hand sides of the equations are polynomials of the
unknown functions). This bijection reduces the problem
to a similar problem in a neighbourhood of the point
(0, 0). Named after I. Bendixson.
References
[I] ANDRONOY, A.A., LEONTOYICH, E.A, GORDON, 1.1. and
MAIER, AG.: Qualitative theory of secondorder dynamic sys
tem~, Wiley, 1973 (translated from the Russian)A FAd
. . n reev
AMS 1980 Subject Classification: 51 N20, 34C20
BENDIXSON TRANSFORMATION - The mapping
4.\
=- -,--,. \' 4\'
= ---
U
x-+v- x'lr'
356
of the Euclidean (x, y)-plane punctured at the point
(0,0) into a similar Euclidean (u, v)-plane. It is the
coordinate expression of the bijection <p generated by
the Bendixson sphere. If the planes (u, v) and (x, y)
coincide, the Bendixson transformation is the inversion
of the plane (x,y) with respect to the circle x 2 +y2 =4.
References
[I] ANDRONOY, A.A., LEONTOYICH, E.A, GORDON, 1.1. and
MAIER, A.G.: Qualitative theory of second order dynamic sys
te~, Wiley, 1973 (translated from the Russian)A FAd
.. n reev
AMS 1980 Subject Classification: 51 N20, 34C20
BERGMAN KERNEL FUNCTION, Bergman kernel -
A function of complex variables with the reproducing
kernel property, defined for any domain D ccn in
which there exist holomorphic functions f=l=0 of class
L 2 (D) with respect to the Lebesgue measure dV. The
function was introduced by S. Bergman [1]. The set of
these functions f forms the Hilbert space
L 2,h(D)CL 2 (D) with orthonormal basis {<PI, </>2, ... };
L 2,h(D)=L 2 (D) n
OeD), where OeD) is the space of
holomorphic functions. The function
z = (z" ... ,zn), ~ = (~" ... '~n)'
is called the Bergman kernel function (or simply the
kernel function) of D. The series on the right-hand side
converges uniformly on compact subsets of D, and
belongs to L 2,h(D) for each given fED, the sum does
not depend on the choice of the orthonormal basis
{<Pj }. The Bergman kernel function depends on 2n com-
plex variables, and is defined III the domain
D X D C C2n ; it has the symmetry property
K(f, z)=K(z, n,it is holomorphic with respect to the
variable z and anti-holomorphic with respect to If r.
D =D' XD", D' cern, D" ccn- rn , then
KD(z, n = K[)(z', nKlJ(z", n
wherez'=(z], ... ,zrn), z"=(zrn+],'" ,zn).
The most important characteristic of the Bergman
kernel function is its reproducing property: For any
function !EL 2,h(D) and for any point z ED the follow-
ing integral representation is valid:
I(z) = jlmK(z, ndVm
[)
Extremal properties of the Bergman kernel function
are:
I) For any point ZED
K(::,::) = sup{i/(z) 1
2 :/r=L 2h iD), 1I/Ik(/J)~l},
2) Let a point fED be such that the class L 2,h(D)
contains functions satisfying the condition fm = 1. The
function K (::, n / K (r, n then satisfies this condition

and has norm K (t f) - I /2, which is minimal for all
such f The function K(z, f) / K(t f) is called the
extremal function of D.
Changes in the Bergman kernel function under biho-
lomorphic mappings are expressed in the following
theorem: If cp is a biholomorphic mapping of a domain
D onto a domain D*, cp(z)=w, CP(f)=1/, then
dz d~
K (
D' W, Tj) = KD(z, n dw dTj'
where dz / dw is the Jacobian of the inverse mapping.
Owing to this property the Hermitian quadratic form
ds 2 -- {-. (l2JogK(z, z) d di
~ ~
J.k = 1 az) (lzk
is invariant under biholomorphic mappings.
The function K(z)=K(z, z), which is also called a
kernel function, plays an important role in the intrinsic
geometry of domains. In the general case it is non-
negative, while the function log K(z) is plurisubhar-
monic. In domains D where K (z) is positive (e.g. in
bounded domains), the functions K(z) and 10gK(z) are
strictly plurisubharmonic. The latter is tantamount to
saying that in such domains D the form ds 2 is positive
definite and. consequently, gives a Hermitian Riemann-
ian metric in D. This metric remains unchanged under
biholomorphic mappings and is called the Bergman
metric. It may be considered as a special case of a
Kahler metric. It follows from extremal property I) that
the coefficients of the Bergman metric increase to infin-
ity on approaching certain boundary points. If D cen
is a strictly pseudo-convex domain or an analytic
polyhedron, then K (z) increases to infinity for any
approach of z to the boundary of the domain. Every
domain which has this property of the Bergman kernel
function is a domain of holomorphy.
For domains of the simplest type, the Bergman ker-
nel function can be explicitly calculated. Thus, for the
ball B = {z: I z I <R} in en, the Bergman function has
]
l
the following form:
n!Rn n_
KB(z, n = -- R2 - ~z)~)
1T" )=1
and for the polydisc U={z: I z) I <Rj , )=1, ... ,n},
in en:
In the special case when n = I and U = B is the disc
{z: I z I <R} in the complex z-plane, the Bergman
metric becomes the classical hyperbolic metric
ds 2 = 2R2 1 dz 12,
(R 2 - 1Z 12)2
which is invariant under conformal mappings and
which defines the LobachevskiI geometry in U.
BERGMAN KERNEL FUNCTION
References
[II BERGMAN, S.: The kernel function and conformal mapping,
Amer. Math. Soc., 1950.
[2] FUKS, B.A.: Special chapters in the theory of analytic functions
of several complex variables, Amer. Math. Soc., 1965 (translated
from the Russian).
[3] SHABAT, B.V.: Introduction to complex analysis, 1-2, Moscow,
1976 (in Russian).
E.M. Chirka
Editorial comments. Recently, additional properties of
the Bergman kernel have been discovered. For a large class
of pseudo-convex domains with Coo boundary, which
includes strictly pseudo-convex domains, and domains of
finite type, the kernel function K(z, w) on 0 X 0 is smooth
up to the boundary in z if w remains fixed in D. This is a
4
consequence of the compactness of the Neumann operator
N for the complex Laplacian on 0 and the identity
P = I-a ''Na.
Here P is the Bergman projection, that is, the orthogonal
projection of L 2 (0) onto L 2h (0) given by integration against
K; a is the Cauchy- Riemann operator and a
its Hilbert
space adjoint. In fact, for these domains P satisfies the so-
called 'condition R', that is P maps L 2s + 2(0) continuously
into L 2s (0), where L 2k (0) denotes the Sobolev space of
order k. This property of the Bergman kernel function is
employed in the study of proper holomorphic and biholo-
morphic mappings. (See [A2], [A4], [AS].) Moreover, the
asymptotic behaviour of K(z, w) has been studied; for
strictly pseudo-convex domains 0 one has
K(z, w) = F(z, w)(il/;(z, w-n-1 +G(z, w) log (il/;(z, w,
where F, G and I/; are Coo functions on 75 X 75 and I/; satisfies
a) I/;(z, z) = p(z) / i, where p is a strictly-plurisubharmonic
defining function for 0;
b) aA and awl/; vanish to infinite order at z= w; and
c) I/;(z, w)= -I/;(w, z).
Similar results have been obtained for certain weakly
pseudo-convex domains, see [A1], [A3], [A4].
The Bergman kernel has also been studied for other
domains, e.g. Cartan domains (cf. [A6]) and Siegel
domains (cl. [A7], Siegel domain).
References
[A 1) BOUTET DE MONVEL, L. and SJOSTRAND, 1.: 'Sur la singular-
-n-I
ite des noyaux de Bergman et de Szego', Asterisque 34-35
(1976), 123-164.
[A2) CATLIN, D.: 'Global regularity of the a-Neumann problem', in
Proc. Symp. Pure Math., Vol. 41, Amer. Math. Soc., 1984,
pp. 39-49.
[A3] DIEDERICH, K., HERBORT, G. and OHSAWA, T.: 'The Bergman
kernel on uniformly extendable pseudo-convex domains',
Math. Ann. 273 (1986),471-478.
[A4) FEFFERMAN, c.: 'The Bergman kernel and biholomorphic
mappings of pseudo-convex domains', Invent. Math. 37
(1974),1-65.
[AS) RANGE, R.M.: Holomorphic functions and integral represen-
tations in several complex variables, Springer, 1986,
Chap!. 7.
[A6) HUA, L.K.: Harmonic analysis of functions of several com-
plex variables in the classical domains, Amer. Math. Soc.,
1963.
[A7] GINDIKlN. S.G.: 'Analysis on homogeneous domains', Rus-
357
BERGMAN KERNEL FUNCTION
sian Math. Surveys 19 (1964),1-89. (Uspekhi Mat. Nauk 19
(1964), 3-92)
AMS 1980 Subject Classification: 32H10
BERGMAN - WElL REPRESENTATION,
Bergman - Weil formula, Wei! formula - An integral
representation of holomorphic functions, obtained by S.
Bergman [I] and A. Weil [2] and defined as follows. Let
D be a domain of holomorphy in en,
let the functions
WI, ... , Uj be holomorphic III D and let
V={zED: I Wk(z) 1<1, k=l, ... ,N} compactly
belong to D. It is then possible to represent ~y func-
tion f holomorphic in V and continuous on V at any
point Z E V by the formula:
f
I fm det (Pi].)
fez) = -'-n ~ n dr,
(2m) II(~.m- ~.(z
k =1
where the summation is performed over all
j 1 < ... <jn, while the integration is carried out over
suitably-oriented n-dimensional surfaces aj" ... ,j,' form-
ing the skeleton of the domain V (d. Analytic
polyhedron), d~=d~11\ ... I\d~n' Here the functions
Pij(~' z) are holomorphic in the domain D X D and are
defined, in accordance with Hefer's lemma [3], by the
equations
n variables', Math. Ann. 111(1935), 178-182.
~m- ~(z) = "2,(ri -z,)Pi/(r, z).
i=1
The integral representation (*) IS called the
Bergman - Wei! representation.
The domains V appearing in the Bergman - Weil
representation are called Wei! domains; an additional
condition must usually be imposed, viz. that the ranks
of the matrices (aJ.fj, /az,,), v=l, ... ,k, fL=l, ... ,n,
k ~n, on the corresponding sets
1= ... =1 ~k I=l}
{ZEV: I~,
are maximal (= k) for all j I < ... <jl: (such Weil
domains are called regular). The Weil dom) ins in the
Bergman - Weil representations may be replaced by
analytic polyhedra U compactly belonging to 0,
U = {zED: Wj(z)ED/, j=1. ... ,N},
where the Dj are bounded domains with piecewise-
smooth boundaries aDj in the plane C. The
Bergman - Weil representation defines the value of a
holomorphic function f inside the analytic polyhedron
U from the values of f on the skeleton a; for n > I the
dimension of a is strictly lower than that of au. If
n = I, analytic polyhedra become degenerate in a
domain with piecewise-smooth boundary, the skeleton
and the boundary become identical, and if. moreover,
N=I and W(z)=z, then the Bergman-Wei! represen-
tation becomes identical with Cauchy's integral for-
mula.
35H
An important property of the Bergman - Weil
representation is that its kernel is holomorphic in z.
Accordingly, if the holomorphic function f is replaced
by an arbitrary function which is integrable over 0,
then the right-hand side of the Weil representation
gives a function which is holomorphic everywhere in U
and almost-everywhere in D \ au; such functions are
called integrals of Bergman - Wei! type. If f is holo-
morphic in U and continuous on fl, then its integral of
Berg~an - Weil type is zero almost-everywhere on
D\U.
Bergman - Weil representations in a Weil domain V
yield, after the substitution
W], (z)
00
~ W,+lm
(*) v ~O lk
the Wei! decomposition
fez) =
="2, ~ Q),j,S,s'{z)(W;,'(z)'" w;:(z
Sir. ~o j I < ... <jk
into a series of functions, holomorphic in D, and this
series is uniformly convergent on compact subsets of V.
References
[I] BERGMAN, S.B.: Mat. Sb. 1 (43) (1936), 242-257.
[2] WElL,A.: 'l'Integrale de Cauchy et des fonctions de plusieurs
[3] VLADIMIROV, V.S.: Methods of the theory of functions of several
complex variables, M.I.T., 1966 (translated from the Russian).
E.M. Chirka
Editorial comments.
References
[A1] HENKIN, G.M. and LEITERl:R, 1.: Theory of functions on com-
plex manifolds, Birkhauser, 1983.
AMS 1980 Subject Classification: 32A25
BERNOULLI AUTOMORPHISM - An automor-
phism of a measure space, which describes Bernoulli
trials and their generalization - a sequence of indepen-
dent trials with the same result and with the same pro-
bability distribution.
Let A be the collection of all possible outcomes of a
trial, and let the probability of the event B CA be given
by the measure v; for a countable set A, denote its ele-
ments by ai and their probabilities by Pi = v(a,). The
phase space of a Bernoulli automorphism is the direct
product of a countable number of copies of the set A.
i.e. the points of the phase space are infinite sequences
h = {hk }. where k runs through the set of integers and
each hk c=A. The transformation T consists in shifting
all members of each sequence one place to the left:
T { hI: } = {hl:,- I }. The measure fl is defi ned as the direct
product of a countable number of measures v; thus if A
is countable. then
 BERNOULLI l.EMNISCATE

In this case, the entropy of the Bernoulli automorphism References

is - V, logp,.
In ergodic theory, Bernoulli automorphisms (or, more
exactly, the cascades generated by iteration of them)
play the role of a standard example of a dynamical sys-
tem, the behaviour of which displays statistical features.
A Bernoulli automorphism is a K-automorphism but
there exist K-automorphisms which are metrically non-
isomorphic to a Bernoulli automorphism, even though
many K-automorphisms are metrically isomorphic to a
Bernoulli automorphism. Two Bernoulli automor-
phisms are metrically isomorphic if and only if they
have the same entropy [1]. A Bernoulli automorphism
is a quotient automorphism of any ergodic automor-
phism of a Lebesgue space with a larger entropy [2].
References
[IA) ORNSTEIN, D.: 'Bernoulli shifts with the same entropy are iso-
morphic', Adv. Math. 4 (1970), 337-352.
[IB) D.ORNSTEIN: 'A Kolmogorov automorphism that is not a Ber-
noulli shift', Matematika 15, no. 1(1971), 131-150 (in Rus-
sian).
[2] SINAi, YA.G.: 'On weak isomorphism of transformations with
invariant measure'. Mat. Sb. 63 (lOS), no. I (1964).23-42 (in
Russian).
D. V Anosov
Editorial comments. A metric isomorphism between two
Bernoulli automorphisms can be given by means of a fini-
tary code ([A1], see also [A4D.
References
[A 1) KEANE, M. and SMORODINSKY, M.: 'Bernoulli schemes with
the same entropy are flnltarily isomorphic', Ann. of Math. 109
(1979), 397-406.
[A2] SMORODINSKY. M.: Ergodic theory, entropy, Springer, 1971.
[A3) ORNSTEIN, D.: Ergodic theory, randomness, and dynamical
systems, Yale Univ. Press, 1974.
[A4) CORNFELD. I.P., FOMIN, S.Y. and SINAI, YA.G.: Ergodic
theory, Springer, 1982 (translated from the Russian).
AMS 1980 Subject Classification: 28DXX
BERNOULLI DISTRIBUTION - A synonym for
binomial distribution.
AMS 1980 Subject Classification: 62EXX
[I] BERNOULLI, J.: Acta Erud. (1695), 59-67; 537-557.
[2] KAMKE, E.: DifJerentialgleichungen: LOsungsmethoden und
LOsungen, 1, Chelsea, reprint. 1971.
N.Kh. Rozov
Editorial comments.
References
[A 1) INCE, W.L.: Ordinary differential equations, Dover, reprint, 1956.
AMS 1980 Subject Classification: 34A08, 34-XX
BERNOULLI INTEGRAL of the equations of hydro-
dynamiCS - An integral which detennines the pressure p
at each point of a stationary flow of an ideal homo-
geneous fluid or a barotropic gas (p =F(p)) in terms of
the velocity v of the flow at that point and the body
force function per unit mass u(x,y, z):
f7=c-iIVI2+U. (1)
The constant C has a specific value for each flow line
and varies from one flow line to another. If the motion
is potential, the constant C is the same for the entire
flow.
For a non-stationary flow, the Bernoulli integral
(sometimes called the Cauchy- Lagrange integral) holds
in the presence of a velocity potential:
f .EP... = ~_ll v 12+u+j(t), (2)
p at 2
where
v = 'V <p(x, y, z, t),
and f (t) is an arbitrary function of time.
For an incompressible liquid the left-hand sides of
equations (l) and (2) are converted to the p / p form;
for a barotropic gas (p =F(p to the form
f.EP... = f p' (p).EP....
p p
The integral was presented by D. Bernoulli in 1738.
References
[ 1] MILNE-THOMSON, L.M.: Theoretical hydrodynamiCS, Macmillan,
1950.
L.N. Sretenskir

BERNOULLI EQUATION - An ordinary first-order AMS 1980 Subject Classification: 76EXX

differential equation BERNOULLI LEMNISCATE - A plane algebraic
ao(x)/ +al(x)y = j(x)yu. curve of order four, the equation of which in orthogo-
where a is a real number other than zero or one. This nal Cartesian coordinates is:
equation was first studied by 1. Bernoulli [1]. The sub- (Xl +y2)2 -2a 2(x 2 _yl) = 0;
stitution i-a = z converts the Bernoulli equation to a
and in polar coordinates
linear inhomogeneous first-order equation, [2]. If a>O,
the solution of the Bernoulli equation is y 0; if p2 = 2a 2 COS 2cp.
O<a< 1, at some points the solution is no longer The Bernoulli lemniscate is symmetric about the coor-
single-valued. Equations of the type dinate origin (Fig.), which is a node with tangents
[f (y)x + g(y )xulr,' = h(y), 0:*0, I, y = +- x and the point of inflection.

~
are also Bernoulli equations if y is considered as the ) .I/~
independent variable, while x is an unknown function
\:

~t;-
ofy.

359
BERNOULLI LEMNISCATE

The product of the distances of any point M to the two
given points FI(-a,O) and F 2 (a, 0) is equal to the
square of the distance between the points F I and F 2.
The Bernoulli lemniscate is a special case of the Cassini
ovals, the lemniscates, and the sinusoidal spirals (cf.
Cassini oval; Sinusoidal spiral).
The Bernoulli spiral was named after lakob Ber-
noulli, who gave its equation in 1694.
References
[1] SAYELOY, A.A.: Plane curves, Moscow, 1960 (in Russian).
D.D. Sokolov
Editorial comments.
References
[A 1] BRIESKORN, E. and KNORRER, H.: Plane algebraic curves,
Birkhauser, 1986 (translated from the German).
AMS 1980 Subject Classification: 53A04
m-1, m-2, ... , and the resulting equations are solved in
xP' x~, ....
If two roots xp and Xq (of multiplicity p and q) of the same
maximum modulus occur, then y(m) satisfies, as m~oo, a
linear difference equation with (x- xp)P(x- Xq)q as charac-
teristic polynomial, and (A1) must be replaced by the
corresponding difference equation of order p + q.
References
[A1] HILDEBRAND, F.: Introduction to numerical analysis,
McGraw-Hili, 1956.
AMS 1980 Subject Classification: 65H05
BERNOULLI NUMBERS - The sequence of rational
numbers B 0, B I, . . . , discovered by 1acob Bernoulli
[1] in connection with the calculation of the sum of
equal powers of natural numbers:

BERNOULLI METHOD - A method for finding the

m -lIn
~ kn= - - ~ s Bsm n + l - s [n)
k =0 n + 1 s=o
real root of algebraic equations of the type
n =0, 1, ... , m = 1,2, ....
(*)
The values of the first Bernoulli numbers are:
with the largest modulus (absolute value). The method
1 1
was proposed by D. Bernoulli [1] and is based on the Bo = 1, BI = -2' B2 = 6' B3 = 0,
following principle. Let y (0), ... ,y (n - 1) be random
1 1
numbers and let the values of y(n),y(n + 1), ... , be B4 = -30' Bs = 0, B6 = 42' B7 = 0,
calculated by the following difference equation:
1 5
aoy(m +n)+aly(m +n -1)+ ... +any(m) = 0. B8 = -30' B9 = 0, BIO = 66' BII = 0.

In general, as m -HY), the expression y (m + 1) / y (m ) All odd-indexed Bernoulli numbers except for B I are
tends to the value of the root of equation (*) with the zero, and the signs of B 2n alternate. Bernoulli numbers
largest modulus. are the values of the Bernoulli polynomials at x =0:
Bn = Bn(O); they also often serve as the coefficients of
References
[1] BERNOULLI, D.: Comm. A cad. Sci. Imper. Petropolitanae 3 the expansions of certain elementary functions into
(1732), 62-69, Petropolis. power series. Thus, for example,
[2] WHITIAKER, E.T. and ROBINSON, G.: Mathematische Bear-
x ox; x'
beitung von Resultaten der Observation, Leningrad-Moscow, -x-
e -1
= ,=0
~Bp-" I x I <2w
1935 (in Russian; translated from the German). P.

L.N. Dovbysh (the so-called generating function of the Bernoulli

Editorial comments. The expression Y(I + 1) / y(m) numbers);
tends to the root of largest modulus only if this root is sim- 00 22,
ple and if no other roots of the same maximum modulus xcotgx = ,~(-I)'B2'(2p)!x2', Ix 1<'17,
occur. If not, then one can proceed as follows (cf. [A1]). Let
_ ox; 22, (22, - 1) 2,
there be one (not necessarily real) root xp of largest tgx - ,~IB2' I (2p)! x, Ix 1<;.
modulus, with multiplicity p. Then the solution of the linear
difference equation for y(m) is L. Euler in 1740 pointed out the connection between
Bernoulli numbers and the values of the Riemann
rp]
y( m) ;::,;
p-l
l
j~O C; j x~, zeta-function rC~) for even s = 2m:

as m--,>x. This expression is the solution of a linear differ- 1'(2 ) - ~2m - (2'17)2m I B I
~ m - :'_( - 2(2n)1 2m'
ence equation with characteristic polynomial (x - xp)P, so
that y(m) satisfies the linear difference equation Bernoulli numbers are used to express many improper
integrals, such as
p-l
y(m) ;::,; ~ xb [p]. y(m- j) (A 1) 1
f '- .-I
x X2n dx I
;=0 J = -4n I B 211 I, n = I, 2.....
II e-"\
as m-HXJ. If p~2, then this equation can be solved expli-
citly in xP ' otherwise (A 1) is rewritten with m replaced by Certain relationships involving Bernoulli numbers are:

360
 BERNOULLI POLYNOMIALS

n'2,l [~]Bk '23 x 2 + '21 x.

B ()_3
3 X - X -
Bo = 1, = 0, n;;;'2,
k =0
Bernoulli polynomials may be computed by the
(the recurrence formula); recurrence formula

n~I[:JBs(X)=nxn-l, n=2,3, ....

s =0

For a natural argument x =m Bernoulli polynomials

The estimates:
2(2n)! < (-1)" -I B ~ ~(2n)!, n ;;;'1,
(21Tfn 2n 3(21T)2n
hold. Extensive tables of Bernoulli numbers are avail-
able; for instance, [2] contains accurate values of B 2n
for n ~90 and approximate values for n ~250.
Bernoulli numbers have found many applications in
mathematical analysis, number theory and approximate
calculations.
References
[I] BERNOULLI, J.: Ars conjectandi, Basle, 1713.
[2] DAVIS. H.T.: Tables of the higher mathematical functions, 2,
Bloomington, 1935.
[3] SAALSCHUETZ, L.: Vorlesungen uber die Bernoullischen Zahlen.
Berlin, 1893.
[4] CHISTYAKOV, I.I.: Bernoulli numbers, Moscow, 1895 (in Rus-
sian).
were first studied by Jacob Bernoulli in 1713, in con-
nection with the computation of the sum
fk n.
k =0
L. Euler [1] was the first to study Bernoulli polynomials
for arbitrary values of x. The term 'Bernoulli polynomi-
als' was introduced by J.L. Raabe in 1851. The funda-
mental property of such polynomials is that they satisfy
the finite-difference equation
Bn(x + 1)-Bn(x) = nx n- I,
and therefore play the same role in finite-difference cal-
culus as do power functions in differential calculus.
Bernoulli polynomials belong to the class of Appell
polynomials, i.e. they satisfy the condition
B~(x) = nBn-l(x)

l
[5] NIELSEN, N.: Traite eiementaire de nombres de Bernoulli, Paris, and are closely connected with the Euler polynomials
1923.
[6] KUDRYAVTSEV, V.A.: Summation of powers of natural numbers
and Bernoulli numbers, Moscow-Leningrad, 1936 (in Russian).
En-I(x) = ~ [Bn(X)-2 nBn [ ] ~
[7] NORLUND, N.E.: Vorlesungen uber DijJerenzenrechnung, Berlin,
1924. The generating function of the Bernoulli polynomials
[8] GELFOND, A.O. [A.O. GEL'FOND]: DijJerenzenrechnung, IS
Deutsch. Verlag Wissenschaft., 1958 (translated from the Rus- ~ = ~ Bn(~) (n.
sian). e t -1 n =0 n.
[9] HARDY, G.H.: Divergent series, Clarendon, 1949.
Bernoulli polynomials are expandable into Fourier
Yu.N. Subbotin
series: For n = 1
Editorial comments. Bernoulli numbers play an impor-
tant role in the theory of cyclotomic fields and Fermat's last BI(x) = x--.l = - ~ sin 27TSX , O<x<l,
2 s= I S7T
theorem, see [A1], pp. 40-41, and [A2]. E.g., if P is an odd
prime number that does not divide the numerators of and for n = 2, 3, ... ,
8 2 , 8 4 , . . . ,8p - 3 , then x P + yP =zP has no solutions in
COS
x, y, ZEN. (See also Cyclotomic field; Fermat great 00

theorem.) Bn(x) = -2n!~--'-------'--, O~x~l.

s= I (21TS)n
References
[A 1] LANG, S.: Cyclotomic fields, Springer, 1978. Bernoulli polynomials satisfy the relationships
[A2] RiBENBOIM, P.: Thirteen lectures on Fermat's last theorem,
Springer, 1979. Bn(mx) = mn-Imi:,IBn [x+...i...]
s=o m
AMS 1980 Subject Classification: 1OL 10, 10H40
(the multiplication theorem);
BERNOULLI POLYNOMIALS - The polynomials Bn(l-x) = (-I)nBn(x)
(the complement theorem);
Bn(x) = s~o [:]Bsnn-s (n=O, 1, ... ),
Bn(x +y) = s~o [:]Bs(y)Xn-s
where Bs are the Bernoulli numbers. Thus, for
n=O,I,2,3, (the theorem on addition of arguments).
Bernoulli polynomials are employed to express the
residual term of the Euler- MacLaurin formula, and for

361
BERNOULLI POLYNOMIALS

the expansion of functions into series. Many important
properties of Bernoulli numbers are a consequence of
the properties of Bernoulli polynomials. Bernoulli poly-
nomials are employed in the integral representation of
differentiable periodic functions
Thus, a Bernoulli random walk may be described in
the following terms. A particle moves 'randomly' along
the x-axis over a lattice of points of the form kh (k is
an integer, h >0). The motion begins at the moment
( = 0, and the location of the particle is noted only at

I/;{x) = 2~ ff(t)dt+
-17
! f !/>k(t)P- I )(x+'1T-t)dt,
-'Tt
discrete moments of time 0, d(, 2dt, .... At each step
the value of the coordinate of the particle increases or
decreases by a magnitude h with probabilities p and
2k -I~ ['1T+t]
!/>k(t) = -k-! -Bk ~ , q = 1 - p, respectively, irrespective of the nature of the
previous motion. Thus, shifts in the positive and in the
and play an important part in the theory of approxima- negative directions ('successes' and 'failures') are
tion of such functions by trigonometric polynomials described by a Bernoulli trial scheme with probability
and other aggregates. Cf. Favard problem. of success p. The Bernoulli random walk is usually
Various generalizations of the Bernoulli polynomials represented geometrically: take the (-axis as the
have been proposed. N.E. Norlund introduced general- abscissa, and the x-axis as the ordinate (cf. Fig. 1,
ized Bernoulli polynomials of order JI and degree n: showing the initial segment of the graph of a random
B~v)(x Iw) = B~v)(x IWI. ... ,wv ) walk of a particle beginning to move from zero).
(certain special cases of such polynomials had been pre-
viously considered by V.G. Immenetskii, N.Ya. Sonin 4h
and D.M. Sintsov). Let 3h
2h
Il rex)
wJ'
= (x +w)- (x)
W h

and D.t i
I 3.6.t 5.c\1 I
B~O)(xlw)=xn, n=O,I, ... ; -h 2.c\1 4.6.1 6.c\t

then BAP)(x Iw) is successively determined as the poly-

nomial solution of degree n of the finite-difference Fig. 1. The initial segment of the graph
equation representing the motion of a particle per-
forming a Bernoulli random walk.
Ilw.B~O)(xlwl"" ,wv) = nB~V_-II)(xlw" ... ,Wv-I),
JI = 1, 2, ... , with initial conditions Let Xj be the random variable corresponding to the
B~P)(Olw" ... ,wp) = B~P)[WI"" ,wv), displacement of the particle in the j-th step. Then
XI, X 2 , , is a sequence of independent random
where B~P)[WI"" ,wp] (a generalized Bernoulli
variables. The coordinate of the randomly-moving par-
number) is found from the recurrence relations
ticle at the moment nM is equal to the sum
~
~
[n]
s WSII B(p)
n-s
= wn"B(v-I)
n-l
Sn = X I + ... + Xn. This is why the graph of a Ber-
s =1 noulli random walk is an illustrative representation of
(B~')[wil = w1Bn , Bd P ) = I, B~O) = 0, Bd P) = I). the behaviour of the cumulative sums of random vari-
ables; moreover, many typical features of the fluctua-
References tions are preserved for the sums of much more general
[1] EULER, L.: Institutiones calculi differentialis, 10, Teubner, 1980.
[2] NORLUND, N.E.: Vorlesungen iiber Differenzenrechnung, Berlin, random variables as well. This graph also shows the
1924. changes in the capital at the disposal of one of the
[3] BATEMAN, H. and ERDELYl, A.: Higher transcendental functions.
players in the classical ruin problem (the formulas giv-
The Gamma function. The hypergeometric function. Legendre
functions, Vol. I, McGraw-Hill, 1953. ing the probabilities of numerous events taking place in
[4] LIKHIN, V.V.: Istor Mat. Issled. (1959),59-134. the course of a Bernoulli random walk were in fact
Yu.N. Subbotin developed in connection with this problem).
AMS 1980 Subject Classification: 33AXX, 1OA 10 A Bernoulli random walk is used in physics as a
rough description of one-dimensional diffusion
BERNOULLI RANDOM WALK - A random walk processes (cf. Diffusion process) and of the Brownian
generated by Bernoulli trials. The example of a Ber- motion of material particles under collisions with
noulli random walk may be used to explain certain molecules.
basic features of more general random walks. In partic- Important facts involved in a Bernoulli random walk
ular, even in this very simple scheme there appear pro- will be described below. In so doing, it is assumed that
perties of 'randomness' which are intuitively paradoxi- 6.t = I, h = 1.
cal. Probabilities of returning. Let the walk begin from

362
 BERNOULlI RANDOM WALK

zero. The probability of at least one return to zero then side more than 9930 units of time with a probability
is I - Ip - q I, that is, one in the symmetric case ;;'0.1, i.e. roughly speaking, such a graph will be
p = q = 1 /2, and less than one if p=!=q. In the sym- observed not less frequently than in one case out of ten
metric case the values of T] (i.e. the time elapsed until (even though this seems absurd at first sight).
the first return to zero) and T2 (the time between the Maximum deviation. If p >q or p <q, the randomly-
first and the second return to zero), etc., are indepen- moving particle will move towards + 00 or - 00 with
dent random variables with an infinite mathematical probability one. Thus, if p<q, one defines the random
expectation. The time elapsed until the N-th return, i.e. variable
the sum T] + ... +TN, increases as N 2 , while the aver- M+ = max S,
O<;;j< 00 J
age number N 2n of returns during 2n steps is given by
and then the probability of M+ =x is
the formula
E(N 2n ) = 2n22n+ I [2n1-
n
I

which increases as Vn: A bounded Bernoulli random walk. One often consid-
Vn
E(N2n)~2y:;;' ers a Bernoulli random walk in the presence of absorb-
ing or reflecting barriers. For instance, let the walk
The consequence is paradoxical: In a symmetric Ber- begin from zero. The presence of an absorbing barrier
noulli random walk, the intervals ('waves') between the at a point a is manifested by that, on reaching this
successive returns to zero on the graph are surprisingly point, the particle ceases to move. In the presence of a
long (Fig. 2).
reflecting barrier at the point a =(k + 1/2), k;;'O an
St 2v'1 integer, the particle will move from k to k - 1 with pro-
800 bability q, and will remain in place with probability p.
600
Finite-difference equations are a principal tool in com-
puting the probabilities of the absorption and of attain-
400 ing specified points. For instance, let the absorbing
barrier be located at the point - a (a >0). If Zt.x is the
200
,,'i probability of a particle located at x at the moment of
o fr---,f-!!J.?:-t<:!Y:-::=---+---~~~:'---_.J.-:-:15,fO~OO:;;"~--
-;-
',.-, ---'Fj2;::O::-:OO::::oo'!- t time t to be absorbed before or at the moment n, then
-\1'''.1 \ { "'I..",; >,,\ ""'; : ,.;,-, the following equation is valid
- 200 I : :' \ _"" . If :t v'....'V~it

- 400
Zt,x = qZt+l,x-1 +pZt+l,x+l, x>-a,
-Vi
with the following obvious boundary conditions:
- 600
Zt, -a = I, O~t~n,
- 800
-2vt Zn,x = 0, x> - a.
Fig. 2. Graphs of three Bernoulli random The solution of this problem for p = q = 1 /2 was
walks: each one was observed during known to A. de Moivre and to P. Laplace. The Laplace
200,000 units of time. formula is:
2 /2 . (
j
1T

Another related feature is that the least probable values = 1-- sm a+x)p( .J..)"d.J.. (*)
of Tn / n (the fraction of time that the graph is above
Zt x
, 'IT .
smcJ>
COS,/, ,/"

the abscissa) are those close to 1/2. More exactly, the where
following theorem is valid: If k---'?oo, n -k---'?oo, then P = a+x+2 [ n-t-x-a]
2 +1.
the formula
p2n,2k ~
'lTn
Vx(l- x) ' The transition to diffusion processes. As an example,
where x = xn,k = k / n, gives the probability p 2n, 2k of let p = q = 1 /2, 6.1 = 1 / N, h = 1 / Vii. Then, as
the equality T 2n = 2k. A corollary is the so-called N ---'?oo, many probabilities, calculated for a Bernoulli
arcsine law: For each O<a< I the probability of the random walk, tend to limits which are equal to the
inequality Tn / n <a tends to respective probabilities of a Brownian motion. As an
example, consider the probability of a particle, which
-Ija vi dx
= -2 .,/
arc Sill Va. has departed from zero, reaching a barrier located at a
'IT x(1-x) 'IT
point 0: before or at the moment T. The limit transition
It is possible to prove, starting from this fact, that dur- from formula (*), for n / N = T, 1 =0, a / Vii =a,
ing 10,000 steps the particle remains on the positive gives the magnitude

363
BERNOULLI RANDOM WALK
1_ _
2_ J e- z'/2dz = _2_ f
a/VT 00
e- z'/2dz,
--Ii,; 0 --Ii,; a/ VT
which is equal to the probability that the coordinate
X(p) of a particle executing a Brownian motion satisfies
the inequality
min X(p) ..;; -a,
O....... T
i.e. to the probability that the particle is absorbed at
the barrier -a. For the sake of a more or less complete
description of such limit relationships it is expedient to
adopt a general standpoint and to consider the transi-
tion from the discrete process of 'cumulative sums' to a
continuous random process (cf. Limit theorems).
The scheme of a Bernoulli random walk is very use-
ful in explaining the characteristic features of the
behaviour of sums of random variables such as the
strong law of large numbers and the law of the iterated
logarithm.
References
[I] FELLER, W.: An introduction to probability theory and its appli-
cations, I, Wiley, 1957-1971, Chapt. 14. Yu. V. Prokhorov
AMS 1980 Subject Classification: 60J15
BERNOULLI SCHEME - A synonym for Bernoulli
trials.
AMS 1980 Subject Classification: 60-XX
BERNOULLI THEOREM - The (historically) original
form of the (weak) law of large numbers. The theorem
appeared in the fourth part of Jacob Bernoulli's book
Ars conjectandi (The art of conjecturing). This part may
be considered as the first serious study ever of proba-
bility theory. The book was published in 1713 by N.
Bernoulli (a nephew of Jacob Bernoulli). The theorem
deals with sequences of independent trials, in each one
of which the probability of occurrence of some event
('success') is p. Let n be the number of trial!. and let m
be the random variable equal to the number of success-
ful events. The Bernoulli theorem states that, whatever
the value of the positive numbers t: and 1/, the probabil-
ity P of the inequality
will be higher than 1-1/ for all sufficiently large n
(n ;;;;:.no). The proof of this theorem, which was given by
Bernoulli and which was exclusively based on a study
of the decrease of probabilities in the binomial distribu-
tion as one moves away from the most probable value,
was accompanied by an inequality which made it possi-
ble to point out a certain bound for the given no if t:
and 1/ were given. Thus, it was found by Bernoulli that
if p =2 / 5, the probability of the inequality
364
__1_ ~ .!!!. _1. ~ _1_
50 -. n 5 -. 50
will be higher than 0.999 if n ;;;;:'25,550. By introducing
a slight improvement in the original reasoning of Ber-
noulli, it is possible to conclude that it is sufficient to
select a value of n obeying the condition
1 + 1 1
n > --Iog-+-,
2 '1/
which gives in turn, for the probability 1- P of the ine-
quality
an estimate of the form
2exp { -+n~}
The condition obtained for the above example is
n ;;;;:'17,665 (more sophisticated estimates show that it is
sufficient to take n ;;;;:'6502; one may note, for the sake
of comparison, that the de Moivre- Laplace theorem
yields 6498 as the approximate value of no). Other esti-
mates for I - P may be obtained using the Bernstein
inequality and its analogues. See also Binomial distribu-
tion.
References
[I] BERNOULLI, J.: Ars conjectandi, Vol. 4, Basle, 1713.
[2] MARKOV, A.A.: Calculus ofprobabilities, Moscow, 1924 (in
Russian).
[3] BERNSHTEiN, S.N.: Probability theory, Moscow-Leningrad, 1946
(in Russian).
Yu. V. Prokhorov
Editorial comments.
References
[A1] FELLER, W.: An introduction to probability theory and its
applications, 1, Wiley, 1965.
[A2] SERFLING, R.I.: Approximation theorems of mathematical
statistics, Wiley, 1980, 6, 96.
AMS 1980 Subject Classification: 60F05
BERNOULLI TRIALS - Independent trials, each one
of which can have only two results ('success' or
'failure') such that the probabilities of the results do
not change from one trial to another. Bernoulli trials
are one of the principal schemes considered in proba-
bility theory.
Let p be the probability of success, let q = 1- P be
the probability of failure, and let 1 denote the
occurrence of success, while 0 denotes the occurrence of
a failure. The probability of a given sequence of suc-
cessful or unsuccessful events, e.g.
1 0 0 1 1 0 1 0 ... I,
is equal to
pqqppqpq "'p=pmqn-m,
where m is the number of successful events in the series
of n trials under consideration. Many frequently occur-
 BERNSTEIN INEQUALITY

ring probability distributions are connected with Ber- [3] KAc, M.: Statistical independence in probability, analysis and
noulli trials. Let Sn be the random variable which is number theory, Math. Assoc. Amer., 1963. A V P kh
. . ro orov
equal to the number of successes in n Bernoulli trials.
AMS 1980 Subject Classification: 60-XX
The probability of the event {Sn = k} is then
BERNSTEIN INEQUALITY - 1) Bernstein's inequal-
[knJ P q ,k -
k n -k -
0, ... ,n,
ity in probability theory is a more precise formulation of
the classical Chebyshev ineqUality in probability theory,
i.e. Sn has a binomial distribution. As n-HYJ, this distri-
proposed by S.N. Bernshtein [1] in 1911; it permits one
bution can be approximated by the normal distribution
to estimate the probability of large deviations by a
or by the Poisson distribution. Let Y 1 be the number of
monotone decreasing exponential function (cf. Probabil-
trials prior to the first success. The probability of the
ity of large deviations). In fact, if the equations
event {Y 1 = k} then is
EX) = 0, EX] = hj , j= I, ... ,n,
qkp, k=O,I, ... ,
hold for the independent random variables XI, ... , Xn
i.e. Y 1 has a geometric distribution. If Yr is the number with
of failures which precede the r-th appearance of a suc-
cessful result, Yr has the so-called negative binomial dis-
tribution. The number of successful outcomes of Ber- (where 1>2 and H is a constant independent of j), then
noulli trials can be represented as the sum the following inequality of Bernstein (where r>O) is
XI + ... + Xn of independent random variables, in valid for the sum Sn=X) + ... +Xn:
which J0 = I if the j-th trial was a success, and Xj =0
otherwise. This is why many important laws of proba- P{ I Sn I >r} ~ 2 exp { - 2(Bnr~ Hr) }, (I)
bility theory dealing with sums of independent vari-
ables were originally established for Bernoulli trial where Bn = ~bj' For identically-distributed bounded
schemes (cf. Bernoulli theorem weak) Law of large random variables Xl (E.Aj =0, EX] =0'2 and I J0 I ~L,
numbers); Strong law of large numbers; Law of the j = 1, ... ,n) inequality (1) takes its simplest form:
iterated logarithm; Central limit theorem; etc.).
A rigorous study of infinite sequences of Bernoulli P{ISnl>t(JVn}~2exp{-2(l::/3)}' (2)
trials requires the introduction of a probability measure
where a=Lt / VnO'. A.N. Kolmogorov gave a lower
in the space of infinite sequences of zeros and ones.
estimate of the probability in (1). The
This may be done directly or by the method illustrated
Bernstein - Kolmogorov estimates are used, in particu-
for the case p = q = 1 /2 below. Let w be a number,
lar, in proving the law of the iterated logarithm. Some
uniformly distributed on the segment (0, 1), and let
idea of the accuracy of (2) may be obtained by compar-
00 X(w) ing it with the approximate value of the left-hand side
w=~~
j =1 2J ' of (2) which is obtained by the central limit theorem in
where J0( w) = 0 or 1, be the expansion of w into a the form
binary fraction. Then the J0, j = 1,2, ... , are indepen-

dent and assume the values and 1 with probability
-2-Je-
V2;
u2 /2du = _2_
I V2;t
[1-~Je-12/2,
t
I / 2 each, i.e. the succession of zeros and ones in the
where 0<0<1. Subsequent to 1967, Bernstein's ine-
expansion of w is described by the Bernoulli trial
qualities were extended to include multi-dimensional
scheme with p = 1 /2. However, the measure on (0, 1)
and infinite-dimensional cases.
can also be specified so as to obtain Bernoulli trials
with any desired p (if r:j= I /2 the measure obtained is References
[I] BERNSHTEiN, S.N.: Probability theory, Moscow-Leningrad, 1946
singular with respect to the Lebesgue measure). (in Russian).
Bernoulli trials are often treated geometrically (cf. [2] KOLMOGOROFF, A.N. [A.N. KOLMOGOROV]: 'Ueber das Gesetz
Bernoulli random walk). Certain probabilities of a large des iterierten Logarithmus', Math. Ann. 101 (1929), 126-135.
[3] HOEFFDING, W.: 'Probability inequalities for sums of indepen-
number of events connected with Bernoulli trials were dent random variables', 1. Amer. Statist. Assoc. 58 (1963), 13-
computed in the initial stage of development of proba- 30.
bility theory in the context of the ruin problem. [4] YURINSKli, V.V.: 'Exponential inequalities for sums of random
vectors', 1. Multivariate Anal. 6 (1976), 473-499.
References A. V Prokhorov
[I] GNEDENKO, B.Y.: The theory of probability, Chelsea, reprint,
1962 (translated from the Russian).
2) Bernstein's inequality for the derivative of a tri-
[2] FELLER, W.: An introduction to probability theory and its appli- gonometric or algebraiC polynomial gives an estimate of
cations, 1-2, Wiley, 1957-1971. this derivative in terms of the polynomial itself. If

365
BERNSTEIN INEQUALITY

Tn(x) is a trigonometric polynomial of degree not BERNSTEIN INTERPOLATION METHOD - A

exceeding n and if sequence of algebraic polynomials converging uniformly
on [ - 1, 1] to a function f (x) that is continuous on this
M = max I Tn{x) I,
O~x~2'1T
interval. More precisely, Bernstein's interpolation
then the following inequalities are valid for all x (cf. method is a sequence of algebraic polynomials
[1]): I T'n'){x) I .;;; Mn', r= 1, 2, ....
n
~A~n)Tn{x)
These estimates cannot be improved, since the number Pn(j;x) = ---,-k_~,--:I_ __ n=I,2, ... ,
Tn(x~n)xx - x~n) ,
M= 1 for
Tn (x ) = cos n{x - xo) where the
Tn(x) = cos{n arc cos x)
is sharp:
max I T'n'){x) I = n'. are the Chebyshev polynomials; the
x

Bernstein's inequality for trigonometric polynomials is

x~n) = cos [ (2k 2~ 1)'/T ]
a special case of the following theorem [2]: If f(x) is
an entire function of order ~a and if are the interpolation nodes; and
M= sup
-oo<x<oo
I/(x)l, A~n) = I(x~n)
then one has if k=l=2/s, I is an arbitrary positive integer, n = 21q + r,
sup Ij<')(x) I .;;; Mo" (r=I,2, .. .). q;;;.l, O~r<2/, s=l, ... ,q; otherwise
-oo<x<oo /-1 /-1
Bernstein's inequality for an algebraic polynomial A?} = ~/(X~?(s-I)+2i+I)- ~/(X~?(s-I)+2i)'
i~O i~1
has the following form [1]: If the polynomial
The ratio between the degree of the polynomial
Pn(x) = ~ akxk Pn(f;x) and the number of points at which Pn(f;x)
k ~o
equals f(x) is (n-l)/(n-q), which tends to
satifies the condition 21 / (2/-1) as n~oo; if I is sufficiently large, this limit
I Pn{x) I .;;; M, a,;;;x';;;b, is arbitrary close to one. The method was introduced
by S.N. Bernstein [S.N. Bernshtein] in 1931 [l].
then its derivative P~(x) has the property
References
I P~{x) I .;;; vex -~~b -x)' a<x<b,
[I] BERNSHTEIN, S.N.: Collected works, Vol. 2, Moscow, 1954, pp.
130-140 (in Russian).
P.P. Korovkin
which cannot be improved. As was noted by S.N.
Bernshtein [1], this inequality is a consequence of the Editorial comments. This method of interpolation seems
not very well known in the West. There is, however, a well-
proof of the Markov inequality given by A.A. Markov.
known method of Bernstein that uses the special interpola-
Bernstein's inequalities are in fact employed in prov-
tion nodes k / n, k=O, ... , n, for bounded functions on
ing converse theorems in the theory of approximation [0, 1]. This method is given by the Bernstein polynomials.
of functions. There are a number of generalizations of The sequence of Bernstein polynomials Bn( f ; x) constructed
Bernstein's inequality, in particular for entire functions for a bounded function f on [0, 1] converges to f(x) at each
in several variables. point of continuity XE[O, 1] of f. If f is continuous on [0, 1],
References the sequence converges uniformly (to f) on [0, 1]. If f is dif-
[I] BERNSTEIN, S.N. [S.N. BERNSHTEIN]: 'Sur l'ordre de la meil- ferentiable, B~( f ; x)-. f(x) (at each point of continuity of
leure approximation des fonctions continues par des f), cf [A 1].
polynomes', Acad. R. Belgique. Cl. Sci. Mem. Coil. 4. Stir. II 4 This method of Bernstein is often used to prove the
(1922).
Weierstrass theorem (on approximation). For a generaliza-
[2] BERNSTEIN, S.N. [S.N. BERNSHTEIN]: 'Sur une propriete des
fonctions entieres', C.R. Acad. Sci. Paris 176 (1923),1603-1605. tion of the method (the monotone-operator theorem), see
[3] NIKOL'SKII, S.M.: Approximation offunctions of several variables [A2], Chapt. 3, Sect. 3. See also Approximation of func-
and imbedding theoremf, Springer, 1975 (translated from the tions, linear methods.
Russian).
N.P. Kornelchuk References
VP. Motornyr [Ai] DAVIS, P.J.: Interpolation and approximation, Dover, reprint,
Editorial comments. 1975.
[A2] CHENEY, E.W.: Introduction to approximation theory, Chel-
References
sea, reprint, 1982.
[Ai] NATANSON, I.P.: ConstructIVe theory ot tunctlOns, 1-3, Ungar,
1964-1965 (translated from the Russian), AMS 1980 Subject Classification: 41 A 10
[A2] LORENTZ, G.G.: Approximation ot tunctlons, Holt, Rinehart
and Winston, 1966.
BERNSTEIN METHOD, method of auxiliary func-
AMS 1980 Subject Classification: 60E15, 42A05, tio/1.1 - A method which is employed in the theory of
41AXX linear and non-linear partial differential equations.

366
 BERNSTEIN METHOD

Bernstein's method consists in introducing certain new

(auxiliary) functions, which depend on the solution +~u [a [~J2
ox
+2b~~+c
ox oy
[~J2]
oy
+
being sought, and which make it possible to establish a
priori estimates of the maximum modulus of the deriva- ou ou u_
+2d-+2e-+g- = Q.
oX oy a
tives of this solution of the required order.
A simple example of the application of Bernstein's Let M be the upper bound of I a I, I b I, I c I in D,
method is the a priori estimate of the modulus of the and let a = 1 /8M. If au / ax and au / ay are con-
derivatives of the solution of the Dirichlet problem for sidered as current coordinates in the plane, and x, y, z
the non-linear (quasi-linear) elliptic equation as parameters, the equation Q =0 is the equation of an
ellipse, since the determinant ale 1 - bt > 2u 2 /3, where
02Z 02 Z
OX 2 + oy2 =f
[ OZ OZ ] _
x,y, Z, ~' oy =
a) = ~ [1+ s~ J, b) = stu' C) = ~ [1+ s~ ].
= [~J2 +2b~~+c [~J2 +
a
ox ox oy oy
(*)
Thus, for any au / ax and au / ay, Q will not be
smaller than a certain negative number - P, Q;:;:' - P
+ 2d~ + 2e~ +g in the disc D, (the number P is readily obtained in explicit form). If
ox oy
one introduces the function u 1 given by the formula
z Ie = 0, P
u) =: U+ 4(x 2 +y2),
where a, b, c, d, e, g are smooth functions of x, y, z;
C is the circle, the boundary of the disc D with radius one obtains
R (the assumption to the effect that D is a disc and 02U) 02U)
z I c = 0 is immaterial, since the general case of an 2 + -2-
-0
x oy
= Q+ P ;;;. 0,
arbitrary-connected domain and inhomogeneous boun-
and u 1 atains its maximum on the boundary C of the
dary condition is readily reduced to the case under con-
domain D and, since u 1 is constant on C, one has
sideration by a change of the function and a conformal
transformation of the domain).
If .fz ;:;:'0, then the estimated maximum modulus n
OU 1
-op ;;;. and ~op""'" 0
:>-, -
1
-PR
2 '
where R is the radius of the circle C. Hence it is possi-
max I z(x,y) I
= (x.I')ED+C ble to find a negative lower bound for az / ap:
n

of the solution of problem (*) is immediately obtained ~ _ ~~:>-, _ aPR

from the maximum principle. op - u op ~ 2e(2n+a)/a'
In order to prove that a regular solution of problem If the same reasoning is applied to a second auxiliary
(*) exists it is sufficient to have a priori estimates of the function u
maximum modulus of the derivatives of the solution up
1
to the third order (cf. Continuation method (to a Z = l/>2(u) =: -n -a+aln - - ,
l-u
parametrized family. To estimate maxc I az / ax I
one obtains an estimate from above
and maxc I az / ay I, it is sufficient to estimate
maxc I az / ap I (since z Ie =0), where p, () are polar ~ ,;::: aP)R (n +a)/2
op '" 2 e
coordinates in the disc D. Now introduce a new (auxili-
ary) function u, given by the formula Thus, maxc I az / ap I is estimated, which means that
maxc I az / ax I and maxc I az / ay I are estimated as
Z = I/(u) = -n -a+alnu, well. The estimate of the maximum modulus of the
where 0:>0 will be selected later. The function first derivatives inside the domain D is performed in a
u =u (x, y) varies from e to e(2n +<>l / <> in the same direc- similar manner: introduce an auxiliary function u given
tion as z(x,y) (-n~z~n). Since by the formula
OZ a OU Z = cJ>3(U) = -n+alnlnu.
ox u ox'
The function u varies in the same direction as z, from e
2n la
to ee . In view of (*), on may write the following
expression for u
and similar results for derivatives with respect to y, it 02U + 02u = _1_ [(l+lnu+aa) [~J2 +2ab~~+
follows that u satisfies the equation ox 2 oy2 u Inu ox ox oy

~+~
ox 2 oy2
= ~U [ [~J2
ox
+ [~J2] +
oy +(1+1nu+ac) [ OUJ2]
oy OU OU ulnu
+2d~+2e oy +g-a- - Q).

367
BERNSTEIN METHOD

Considerations similar to those given above show that segment O~x ~ I, the convergence
if the function B~kV;X) ~ j<k)(X)

is uniform on this segment. A study was made ([IB],

attains a maximum in the domain D, this maximum
does not exceed some number, the value of which
depends solely on nand M. This yields the required
estimates of maxD 1 oz / ox
1 and maxD 1 I. oz / oy
Bernstein's method may also be used to estimate, in
a similar manner, the maximum modulus in the domain
D + C of all highest derivatives of the solution (the only
other operation which is required is the differentiation
of the initial equation).
The method was first utilized by S.N. Bernstein [S.N.
Bernshteln] [I]. The method was subsequently
extended and was systematically utilized in the study of
various problems for elliptic and parabolic differential
operators [3], [4], [5].
References
[IA] BERNSTEIN, S.N. [S.N. BERNSHTEIN]: 'Sur la generalisation du
probleme de Dirichlet (premiere partie)', Math. Ann. 62 (\ 906),
253-271.
[lB] BERNSTEIN, S.N. [S.N. BERNSHTEIN]: 'Sur la generalisation du
probleme de Dirichlet (deuxieme partie)', Math. Ann. 69
(\ 910), 82-136.
[5]) of the convergence of Bernstein polynomials in the
complex plane if f is analytic on the segment O~x ~ 1.
References
[IA] BERNSHTEiN, S.N.: Collected works, Vol. I, Moscow, 1952, pp.
105-106.
[IB) BERNSHTEiN, S.N.: Collected works, Vol. 2, Moscow, 1954, pp.
310-348.
[2) GoNCHAROY, V.L.: The theory of interpolation and approxima-
tion offunctions, Moscow, 1954 (in Russian).
[3) BASKAKOV, V.A.: 'An instance of a sequence of linear positive
operators in the space of continuous functions', Dokl. Akad.
Nauk SSSR 113, no. 2 (1957), 249-251 (in Russian).
[4) KOROVKIN, P.P.: Linear operators and approximation theory,
Hindushtan Publ. Comp., Delhi, 1960 (translated from the Rus-
sian).
[5) KANTOROYICH, L.V.: Izv. Akad. Nauk SSSR Ser. Mat. 8
(1931), 1103-1115.
P.P. Korovkin
Editorial comments. There is also a multi-variable gen-
eralization: generalized Bernstein polynomials defined by the
completely analogous formula
BnU, Xl,'" ,Xk) =
~
n,

i, =0
~ f[~'
1,=0nl
... ,i...] [~l]
nk 1
[~k]k X
[2] BERNSHTEIN, S.N.: Collected works, Vol. 3, Moscow, 1960.
XX;'(1-Xlt,-I,. X~(1-Xk)n'-I'.
[3] LADYZHENSKAYA, O.A. and URAL'TSEYA, N.N.: Linear and
quasilinear elliptiC equations, Academic Press, 1968 (translated Here n stands for the multi-index n = (nl , ... ,nk)'
from the Russian). As in the one variable case these provide explicit polyno-
[4] POGORELOY, A.V.: Die Verbiegung konvexer Fliichen,
Akademie-Verlag, 1957 (translated from the Russian). mial approximants for the more-variable Weierstrass approx-
[5] OLEJNIK, O.A. and KRUZHKOY, S.N.: 'Quasi-linear parabolic imation and Stone-Weierstrass theorems. For the
equations of second order in several independent variables', behaviour of Bernstein polynomials in the complex plane
Russian Math. Surveys 16, no. 2 (1961), 105-146. (Uspekhi and applications to movement problems, cf. also [A3].
Mal. Nauk 16, no. 5 (\961),115-155)
I.A. Shishmarev References
[A1) DAYIS, P.J.: Interpolation and approximation, Dover, reprint,
AMS 1980 Subject Classification: 35J65, 35K60
1975, pp. 108-126.
[A2] RIVLIN, TH.J.: An introduction to the approximation of func-
BERNSTEIN POLYNOMIALS - Algebraic polynomi- tions, Dover, reprint, 1981.
als defined by the formula [A3) LORENTZ, G.G.: Bernstein polynomials, Univ. of Toronto
Press, 1953.
Bnif; x) = Bn(x) =
AMS 1980 Subject Classification: 33AXX, 41 A 10,
k~![ ~ j [~jxk(l_x)n-k, n=1,2, .... 41A17

Introduced by S.N. Bernshtein in 1912 (cf. [1]). The BERNSTEIN - ROGOSINSKl SUMMATION
sequence of Bernstein polynomials converges uniformly METHOD - One of the methods for summing Fourier
to a function f on the segment O~x ~ 1 if f is continu- series; denoted by (BR, an). A trigonometric series
ous on this segment. For a function which is bounded
by C, 0 < C < 1, with a discontinuity of the first kind,
f( C - ) +f( C + ) is summable by the Bernstein - Rogosinski method at a
Bnif; C) ~ 2 .
point Xo to the value S if the following condition IS
The equation satisfied:

Bnif; c)- f(c) = (c)c(l-c)

2n
+0 [11
'n

is valid if f is twice differentiable at the point c. If the

k-th derivative f(k) of the function is continuous on the

36R
 BERRY-ESSEEN INEQUALITY

where {an}, an >0, an~O, is a sequence of numbers, minimal surfaces; if k <n - 1, then already for n;;;:.4 it
and where the Sn(x) are the partial sums of the series is possible to find non-linear minimal surfaces Fk,
(*). defined over any Ek.
W. Rogosinski [1] first (1924) considered the case References
an = p7T / 2n, where p is an odd number, and then [I] BERNSTEIN, S.N. [S.N. BERNSHTEiN]: 'D'ber ein geometrisches
(1925) the general case. S.N. Bernstein [S.N. Theorem und seine Anwendung auf die partiellen Differential-
gleichungen vom elliptischen Typus', Math. Z. 26 (1927),551-
Bernshtein] [2] considered (1930) the case 558 (translated from the French).
an =7T / (2n + 1). The (BR, an)-method sums the Fourier [2] NITSCHE, J.e.e.: Vorlesungen uber Minimaljldchen, Springer,
series of a function f EL [0, 27T] in the cases an =p7T / 2n 1975.
[3] OSSERMAN, R.: 'Minimal varieties', Bull. A mer. Math. Soc. 75
and an =7T / (2n + 1) at the points of continuity of the (1969), 1092-1120.
function to its value and is one of the regular :summa- [4] OSSERMAN, R.: A survey of minimal surfaces, v. Nostrand, 1969.
tion methods. [5] FOMENKo, A.T.: Plateau's problem, Gordon and Breach, 1987
(translated from the Russian).
The Bernstein - Rogosinski sums Bn(x, an) are I.Kh. Sabitov
employed as an approximation procedure. In both
Editorial comments. As an important reference for the
cases described above they realize an approximation of generalizations of Bernstein's theorem the paper of
the same order as the best approximation for functions Bombieri~de Giorgi~Giusti [A1] can be Quoted. The origi-
of the classes Lip a and WI Lip a. nal of [1] is [A2].
References References
[I] ROGOSINSKl, W.W.: 'Ueber die Abschnitte trigonometischer [A 1] BOMBIERI, E., GIORGI. E. DE and GIUSTI, E.: 'Minimal cones
Reihen', Math. Ann. 95 (1925), 110-134. and the Bernstein theorem', InventlOnes Math. 7 (1969),
[2] BERNSHTEiN, S.N.: Collected works, Vol. I, Moscow, 1952, p. 243-269.
37 (in Russian). [A2] BERNSTEIN, S.N.: 'Sur une theoreme de geometrie et ses
[3] HARDY, G.H.: Divergent series, Clarendon, 1949. applications aux equations derivees partlelles du type ellip-
A.A. Zakharov tique', Comm Soc. Math. Kharkov 15 (1915-1917), 38-45.
Editorial comments. AMS 1980 Subject Classification: 53A 10, 53A07
References
[A1] BEEKMANN, W. and ZELLER, K.: Theorie der Limitierungsver- BERRY - ESSEEN INEQUALITY - An inequality
fahren, Springer, 1970. giving an estimate of the deviation of the distribution
AMS 1980 Subject Classification: 40GXX, 42A24 function of a sum of independent random variables
from the normal distribution function. Let X I, . . . ,Xn
BERNSTEIN THEOREM on minimal surfaces - If a be independent random variables with the same distri-
minimal surface is given by the equation z =/ (x, y), bution such that '
where / has continuous partial derivatives of the first
EX} = 0, EX] = 02 > 0, E 1 JS 13 < 00.
and second orders for all real x and y, then F is a
plane. A proof of this theorem, which is due to S.N. Let
Bernstein [S.N. Bernshtein] [1], is a consequence of a
more general theorem on the behaviour of surfaces with and
= -- f
1 x 1
non-pOSitive curvature. Various generalizations of <l>(x) e- r l2dt;
Bernstein's theorem have been proposed, most of them yf2; -00

being of the three following kinds: 1 Quantitative then, for any n,

,I; X}';;;;X}-<l>(X) I ,; ; A+,

improvements; e.g. obtaining a priori estimates of the
form I K(O) I ~const/ R2 where R is the radius of the
disc over which the minimal surface z = / (x, y) is
defined and K(O) is the Gaussian curvature of the sur-
face at the centre of the disc. 2) The search for other a
priori geometric conditions under which the minimal
surface would be of a specific kind - a plane, a
catenoid, etc.; for instance, if the spherical image of a
complete minimal surface contains no open set on the
sphere, then such a minimal surface is a plane. 3) The
generalization of Bernstein's theorem to minimal sur-
faces Fk of dimension k, located in a Euclidean space
En; for example, if k = n -1, any minimal surface over
all En - I is uniquely determined if n ~ 8, and is a
hyperplane, while if n >8, there exist non-planar
sup Ip{ ovn;=1
x vn
where A is an absolute positive constant. This result
was obtained by A.e. Berry [1] and, independently, by
e.G. Esseen [2].
References
[I] BERRY, A.e.: 'The accuracy of the Gaussian approximation to
the sum of independent variables', Trans. Amer. Math. Soc. 49
(1941), 122-136.
[2] EsSEEN, e.G.: 'On the Liapunoff limit of error in the theory of
probability', Ark. Mat. Astr. Fysik 28A, no. 2 (1942), 1-19.
[3] I'ETROV, V.V.: Sums of independent random variables, Springer,
1975 (translated from the Russian).
v. V. Petrov
Editorial comments. The constant A can be taken to be
33/4, cf. [A1], p. 515 fl.

369
BERRY - ESSEEN INEQUALITY
References
[A 1] FELLER, W.: An introduction to probability theory and its
applications, 2, Wiley, 1966.
AMS 1980 Subject Classification: 60E15
BERTINI THEOREMS - Two theorems concermng
the properties of linear systems (cL Linear system) on
algebraic varieties, due to E. Bertini [1].
Let V be an algebraic variety over an algebraically
closed field k of characteristic 0, let L be a linear sys-
tem without fixed components on V and let W be the
image of the variety V under the mapping given by L.
The following two theorems are known as the first and
the second Bertini theorem, respectively.
1) If dim W> 1, then almost all the divisors of the
linear system L (i.e. all except a closed subset in the
parameter space P (L) not equal to P (L are irreduci-
ble reduced algebraic varieties.
2) Almost all divisors of L have no singular points
outside the basis points of the linear system L and the
singular points of the variety V.
Both Bertini theorems are invalid if the characteristic
of the field is non-zero.
Conditions under which Bertini's theorems are valid
for the case of a finite characteristic of the field have
been studied [3], [6]. If dim W= 1, Bertini's theorem is
replaced by the following theorem: Almost all fibres of
the mapping CPL: V ~ Ware irreducible and reduced if
the function field k (W) is algebraically closed in the
field k(V) under the imbedding cP~: k(W)~k(V). If
the characteristic of k is finite, the corresponding
theorem is true if the extension k (V) / k (W) is separ-
able [3], [6]. The Bertini theorems apply to linear sys-
tems of hyperplane sections, without restrictions on the
characteristic of the field [5].
References
[I] BERTINI. E.: Introduction to the projective geometry of hyper-
spaces. Messina, 1923 (in Italian).
[2] Algebraic surfaces. Moscow. 1965 (in Russian).
[3] BALDASSARRI, M.: Algebraic varieties, Springer, 1956.
[4] AKIZUKI. Y.: 'Theorems of Bertini on linear systems', 1. Marh.
Soc. Japan 3, no. 1(\951).170-180.
[5] NAKAI, Y.: 'Note on the intersection of an algebraic variety
with the generic hyperplane'. Mem. Call. Sci. Univ. K\,oto Ser.
A Marh. 26. no. 2 (1950). 185-187.
[6] ZARISKI. 0.: 'The theorem of Bertini on the variable singular
points of a linear system of varieties'. Trans. Amer. Math. Soc.
56 (1944). 130-140.
[7] HARTSHORNE. R.: Algebraic Keometrv. Springer. 1977.
V.A. Iskovskikh
AMS 1980 Subject Classification 14C20
BERTRAND CRITERION of convergence of series
' " ryJ
~n 1
a fl with f70sitive numhers as terms - If .
= rn l' ~ -
E"
l a" I I
I j - Illn 11
J three different problems are in fact solved. If the distri-
370
and if the limit (finite or infinite)
B = lim Bn
n~oc
exists, then the series is convergent if B> 1 and IS
divergent if B < 1. Established by J. Bertrand.
References
[I] FICHTENHOLZ, G.M.: Differential und Integralrechnung. I,
Deutsch. Verlag Wissenschaft., 1964.
L.D. Kudryavtsev
AMS 1980 Subject Classification: 40A05
BERTRAND CURVES, conjugate curves, Bertrand
pair - Two space curves L and L with common princi-
pal normals. Let k I and k 2 be the curvature and the
torsion of L respectively. For the curves L and L to
be conjugate it is necessary and sufficient that
ak J sinw+ak 2 cosw = sinw
is true. Here a is a constant, and w is the angle between
the tangent vectors of L and L. The name Bertrand
curve is also given to a curve L for which there exists a
conjugate curve L. They were introduced by J. Ber-
trand in 1850.
E. V. Shikin
Editorial comments. Bertrand's original paper is [A2]. A
general reference is [A 1].
References
[A 1] BLASCHKE, W. and LEICHTWEISS, K.: Elementare Differential-
geometrie, Springer, 1973.
[A2] BERTRAND, J.: 'Memoire sur la theorie des courbes a double
courbure', Liouvilles Journal 15 (1850).
AMS 1980 Subject Classification: 53A04
BERTRAND PARADOX (in probability theory) - A
paradox connected with an inaccurate formulation of
the initial assumptions in solving problems in probabil-
ity. Noted by J. Bertrand [1]. Bertrand's problem is con-
cerned with the probability that the length of a chord,
chosen at random in a disc of radius one, is larger than
the side length of the inscribed equilateral triangle. Ber-
trand singled out three various values of the unknown
probability (1 /2, 1 /3, I /4), depending on the
parameters characterizing the location of the chord. (In
the first case by the distance p to the centre of the cir-
cle and the angle (j between the normal to the chord
and the x-axis; in the second by the angular coordi-
nates ex and f3 of the points of intersection of the chord
with the circumference of the disc; in the third case by
the Cartesian coordinates (x, y) of the base of the nor-
mal dropped from the centre of the disc. In all three
cases the coordinate origin coincides with the centre of
the disc.) It was shown by H. Poincare [2] that the ori-
gin of the paradox is due to the fact that in each of the
three cases the respective pair of parameters is assumed
to be uniformly distrihuted in the given zone, so that
 BESSEL EQUATION

bution of some pair (say, IX and {3) is fixed, the distri- sumrnable to degree p on any finite interval of the real
bution of all other parameters can be uniquely calcu- axis, is called a Besicovitch almost-periodic function if
lated (and is not necessarily uniform, even if IX and f3 to each >0 there corresponds a sufficiently homogene-
are uniformly distributed). From the geometrical point ous set of numbers (the so-called (BP, )-almost-periods
of view, the most natural assumption is that p and () of f(x)):
are independent and uniformly distributed in the inter- ... < 'T -2 < 'T -I < 'To < 'TI < ... ,
val O~p~ 1, 0~()~2'lT (cf. [3]).
such that for each i
References
[I) BERTRAND, J.L.: Calcul des probabilites, Gauthier-Villars, 1907. Mx{ I I(x +'Ti)-I(x) IP} < f!',
[2) POINCARE, H.: Calcul des probabilites, Gauthier-Villars, 1912.
[3) KENDALL, M.G. and MORAN, P.AP.: Geometric probability, and for each c >0
Griffin, 1963.
A. V Prokhorov
Editorial comments.

f
References
where
[A 1) SZEKELY, G.G.: Paradoxes in probability theory and 1 F(x)dx,
Mx{F(x)} = lim -2
7----+00 T ~T
mathematical statistics, Reidel, 1986, pp. 43-48.
- - 1 n
AMS 1980 Subject Classification: 60005, 62AXX M,{F(i)} = n~~ 2n + 1 i~nF(i).
BERTRAND POSTULATE - For any natural number Here F(x) is a real-valued function, defined, respec-
n > 3 there exists a prime number that is larger than n tively, for a real variable and an integer argument.
and smaller than 2n - 2. In its weaker formulation
References
Bertrand's postulate states that for any x> 1 there [1) BESICOVITCH, AS.: 'On mean values of functions of a complex
exists a prime number in the interval (x, 2x). The pos- and of a real variable', Proc. London Math. Soc. (2) 27 (1927),
tulate was advanced by J. Bertrand in 1845 on the 373-388.
[2) BESICOVITCH, AS.: 'On Parseval's theorem for Dirichlet series',
strength of tabulated data. and was proved by P.L Proc. London Math. Soc. (2) 26 (1927),25-34.
Chebyshev (cf. Chebyshev theorems on prime numbers). [3) LEVITAN, B.M.: Almost-periodic junctions, Moscow, 1953 (in
Russian).
References E.A. Bredikhina
[I) CHEBYSHEV, P.L: Oeuvres de P.L. Chebyshev, I, Chelsea.
reprint, 1961 (translated from the Russian). B M B d'kh' Editorial comments. Besicovitch developed his theory in
. . re I In
[A1], rather than in [1], [2].
Editorial comments. As is implicit in the article, for each p;;;.1 there is a class
of almost-periodic functions, denoted by 8 P . The first part of
References
[A 1) HARDY. G.H. and WRIGHT, E.M.: An introduction to the
the article deals with 8 2 , the rest is more general. General
theory of numbers, Clarendon Press, 1965, p. 343ff. references may be found under Almost-periodic function.
References
AMS 1980 Subject Classification: 1OL 10, 10H15 [A 1) BESICOVITCH, AS.: 'On generalized almost periodic func-
tions', Proc. London Math. Soc. (2) 25 (1926),495-512.
BESICOVITCH ALMOST-PERIODIC FUNCflONS -
A class of almost-periodic functions in which the analo- AMS 1980 Subject Classification: 42A75
gue of the Riesz- Fischer theorem is valid: Any tri-
gonometric series BESSEL EQUATION - A second-order linear ordi-
nary differential equation
x 2/ ' + xi +(x 2 -'?)y = 0, p=const, (I)
where
or, in self-adjoint form:

is the Fourier series of some B 2 -almost-periodic

func-
(xy ") + [x - ,? ]y
~ = o.
tion. The definition of these functions [I], [2] is based
on a generalization of the concept of an almost-period, The number JI is called the order of the Bessel equation;
and certain additional ideas must be introduced in it. in the general case x, y and JI assume complex values.
A set E of real numbers is called sufficiently homo- The substitution y = ux -I /2 yields the reduced form of
geneous if there exists an L >0 such that the ratio equation (1):
between the largest number of members of E in an
interval of length L and the smallest number of u"+ [ I+~
1-4'? ] u = O. (2)
members in an interval of the same length L is less
than 2. A sufficiently homogeneous set is also relatively A Bessel equation is a special case of a confluent
dense. A complex-valued function f (x), - 00 <x < 00, hypergeometric equation; if x = z /2i is substituted into

371
BESSEL EQUATION

(2), equation (2) becomes a Whittaker equation. In
equation (1) the point x =0 is weakly singular, while
the point x = 00 is strongly singular. For this reason a
Bessel equation does not belong to the class of Fuchs-
ian equations (cf. Fuchsian equation). F. Bessel [1] was
the first to study equation (1) systematically, but such
equations are encountered even earlier in the works of
D. Bernoulli, L. Euler and J.L. Lagrange.
A Bessel equation results from separation of vari-
ables in many problems of mathematical physics [2],
particularly in the case of boundary value problems of
potential theory for a cylindrical domain.
The solutions of Bessel equations are called cylinder
functions (or Besselfunctions). These may be subdivided
into the cylinder functions of the first kind (Bessel
functions) J.(x), the cylinder functions of the second
kind (Weber functions or Neumann functions, (cf.
Weber function; Neumann function) Y.(x) and the
cylinder functions of the third kind (Hankel functions)
m
H\}l (x), 2) (x). If the order p is fixed, all these func-
tions are analytic functions of the complex argument x;
for all these functions, except for the functions In(x) of
integer order, the point x = 0 is a branch point. If the
argument x is fixed, all these functions are single-
valued entire functions of the complex order P [3].
If the order p is not an integer, then the general solu-
tion of equation (1) may be written as
y = C 1J v(X)+C 2J -vex),
where C 1, C 2 are arbitrary constants. For a given
order, any two of the functions J .(x), Y .(x), HVl (x),
The substitution y =x'w transforms (1) into the
Laplace equation:
xw" +(21'+ l)w' +xw = 0;
which permits one to represent the solutions of (1) by
contour integrals in the complex plane.
In applications it is often required to find the eigen
values of the equation
x 2y" +xy' +(Ax2_~)y = 0, (3)
where p is fixed while A is a parameter. Equation (3) on
the segment O:S;;;;x';;;;;a with the boundary conditions:
y(x) is bounded as x~o, yea) = 0,
is an example of a problem with a discrete spectrum
(the e1gen values are determined by the condition
J .(a VA )=0 in terms of the zeros of a Bessel function).
Equation (3) with the boundary condition:
y(x) is bounded on the semi-axis O.;;;x<oo,
represents a problem with a continuous spectrum (eigen
values A~O).
The inhomogeneous Bessel equation
x 2y" +xy' +(X2_~)y = f(x) (4)
has the particular solution
y =; Yv(x)jxJv(x)f(x)dx- ; Jv(x)jxYv(x)f(x)dx.
Solutions of equation (4) have been studied in more
detail for a right-hand side of special form. Thus, if
f(x)=x P, equation (4) is satisfied by a Lommel func-
tion; if 4( v+ I

m 2 ) (x) are linearly independent and may serve as a

f(x) = _x /2)
y;f(p+ 1/2)
fundamental system of solutions of (1). For this reason, it is satisfied by a Struve function; if
the general solution of equation (1) can be represented,
in particular, in the following forms: f(x) = -.L(x-p)sinJI'IT,
'IT
y = C 1J v(X)+C 2 Yv(x), y = C 1HV)(x)+C 2 He)(x). it is satisfied by an Anger function; and if
The following equations are closely connected with
1
equation (l): the equation f(x) = --[(x+p)+(X-V)COSJI'IT],
'IT

z2y " +zy' _(Z2 +~)y = 0, it is satisfied by a Weber function.

which becomes (1) as a result of the substitution z = ix, There are linear equations of higher orders with solu-
and with as a fundamental system of solutions the tions whose properties are analogous to those of Bessel
modified cylinder functions (Bessel functions of imaginary functions. The general n-th order equation of Bessel type
argument), and the equation has the form
Z2 y " +zy' _(iz2+~)y = 0,
which becomes equation (1) as a result of the substitu-
IT
,=I
[x~
dx
+c,lJy+xny = 0,
tion z = Vix and which has the Kelvin functions as its
fundamental system of solutions. Many other second- C, = const, ~CI = 0,
order linear ordinary differential equations (e.g. the
Airy equation) can also be transformed into equation and its solution depends on n - 1 parameters. In partic-
(1) by a transformation of the unknown function and ular, a third-order equation of Bessel type (which has a
the independent variable. The solution of a series of solution with two parameters lX, {3) may be represented
linear equations of higher orders may be written in the in the form:
form of Bessel functions (4). xY" +3x 1y" +[1 +9a,B-3(a+f3)lJxy' +

372
 BESSEL INEQUALITY

+[x 3-9ap(a+p)+2(a+p)3]y = 0, conditions the following expansion is valid:

a, P = const _ 00
f(x) - ~cnJp - , xJ'
[JL~ 1
.-1
References
[ JL~
[1) BESSEL, F.: Abh. d. K. Akad. Wiss. Berlin (1824), 1-52.
(2) GRAY, A. and MATHEws, G.B.: A treatise on Besselfunctions c. = ,2J~+12 (JL~) [f(x)Jp
I ]
-,-x x dx, o<x<,.
and their application to physics, Macmillan, 1931.
(3) WATSON, G.N.: A treatise on the theory of Bessel functions, 1-2, In an infinite interval this expansion is replaced by the
Cambridge Univ. Press, 1952.
Fourier- Bessel integral
(4) KAMKE, E.: Differentialgleichungen: Losungsmethoden und
Losungen, 1, Chelsea, reprint, 1971. 00 00
N.Kh. Rozov f(x) = [c~..Jp(Ax)dA, CA = [f(X)Jp(Ax)xdx,
Editorial comments.
o<x<oo.
References
[A1) LEBEDEV, N.N.: Special functions and their applications, The following formulas play an important role in the
Dover, reprint, 1972 (translated from the Russian).
theory of Bessel functions and their applications:
AMS 1980 Subject Classification: 34830, 33A40 1) the integral representation

BESSEL FUNCflONS - Cylinder functions of the In(z) = -I "

fcos(z sin<p-n<p)d<p,
'1T 0
first kind. A Bessel function of order p can be defined
as the series 2) the generating function
_ 00 (-It [Z]P+2k (*) ez(~-Cl)/2 = ~ In(z)/!,,
Jp(z) - k~or(k+l)r(k+p+l) "2 n= -00

3) the addition theorem for Bessel functions of order

_ [z]P (-It [z]2k
- "2 k~or(k+l)r(k+p+l) "2 '
00
zero
Jo(Va2+b2-2abcos<p) =
which converges throughout the plane. A Bessel func- 00

tion of order p is the solution of the corresponding

= J o(a)Jo(b)+2 ~Jk(a)Jk(b)cosk<P,
k=1
Bessel equation. If the argument and the order v are
4) the recurrence formulas
real numbers, the Bessel function is real, and its graph
has the form of a damped vibration (Fig.); if the order Jp_I(Z)+Jp+l(z) -2:P..
- Jp(z),
z
is even, the Bessel function is even, if odd, it is odd.
[JpCz)r = t[Jp-I(Z)-Jp+l(z).
For references, see Cylinder functions.
P.I. Lizorkin
I
AMS 1980 Subject Classification: 33A40

Graphs of the functions y =J 0 (x) and BESSEL INEQUALITY - The inequality

y=J 1(x).
The behaviour of a Bessel function in a neighbourhood
of zero is given by the first term of the series (*); for
large x, the asymptotic representation
Jp(x) ~ .,.V12cos
-:;;; [x-:!!...v-:!!...]
2 4
holds. The zeros of a Bessel function (i.e. the roots of
the equation J .(x) =0) are simple, and the zeros of
J .(x) are situated between the zeros of J.+ 1(x). Bessel
functions of 'half-integral' order v=n + 1/2 are expres-
sible by trigonometric functions; in particular
IIfl12 = (j,j);;;;' ~ I (j, <Pa) 12
aEA (<Pa, <Pa)
a~ IVII:: II ] 12,
where f is an element of a (pre-) Hilbert space H with
scalar product (j, </ and {</>a: aEA} is an orthogonal
system of non-zero elements of H. The right-hand side
of Bessel's inequality never contains more than a count-
able set of non-zero components, whatever the cardinal-
ity of the index set A. Bessel's inequality follows from
the Bessel identity

J I/ 2(X) = .,.
V12
-:;;; sin x, J -1/2(X) = .,.
Vf2
-:;;; cosx. Ik-i~xa;<Pa, Ir IfI2-i~Aa; I x a, 12,
The Bessel functions J .(p.~/-1 x) (where p.~ are the posi- which is valid for any finite system of elements {</>a,:
tive zeros of J .(x), v> -1 /2) form an orthogonal sys- i = 1, ... ,n}. In this formula the x a; are the Fourier
tem with weight x in the interval (0, I). Under certain coefficients of the vector f with respect to the orthogo-

373
BESSEL INEQUALITY

nal system {cf>a l , ,cf>a.}, i.e. xo, Xo +h, Xo -h, ... ,Xo +nh, Xo -nh, Xo +(n + I)h,
at the point x=xo+th:
xa, = fif' <Pa), Aa; = (<Pa;, <PaJ
G 2n +2(XO+th) =/0+1112t+1o ~
a;
2! + ... + (1)
The geometric meaning of Bessel's inequality is that
the orthogonal projection of an element f on the linear + n+lt(t2-I),(t 2 -n 2)
fil2 (2n+I)! '
span of the elements cf>a, aEA, has a norm which does
not exceed the norm of f (i.e. the hypothenuse in a and the Gauss formula of the same order for backward
right-angled triangle is not shorter than one of the interpolation with respect to the node Xl =xo+h, i.e.
other sides). For a vector f to belong to the closed with respect to the population of nodes
linear span of the vectors cf>a, aEA, it is necessary and Xo +h, Xo, Xu +2h, Xu -h, ... ,Xo +(n + I)h, Xo -nh:
sufficient that Bessel's inequality becomes an equality. t{t - }\
If this occurs for any fEH, one says that the Parseval G 2n +2(XO+th) = /1 +1112(t-I)+fi~+ ... + (2)
2!
equality holds for the system {cf>a: aEA} in H. + n+1 t(t 2 -I)'" [t 2 -(n-1)2](t-n)(t-n-I)
Fora system {cf>a: a = 1, 2, ... } of linearly indepen- fij2 (2n+I)! .
dent (not necessarily orthogonal) elements of H Bessel's Putting
identity and Bessel's inequality assume the form

Ik-a,~ b~Pif, I <Pp)<Pa /1

2
Bessel's interpolation formula assumes the form ([ 1],
[2]):
B2n+2(XO+th) = (3)

= /1/2 +11/2 [t - tJ+ fi I 2 t(t; I) + ... +

II! 1/ 2 ;;;. ~ b~Pif, <Pa)(/, <Pp),
a,p=1
+ n t(t 2 -1)'" [t 2 -(n-I?](t-n) +
where b~{J are the elements of the matrix inverse to the fil2 (2n)!
Gram matrix (d. Gram determinant) of the first n vec- + n + I t(t 2 -1) ... [t 2 -en -1)2](t -n)(t - I ! 2)
tors of the initial system. fi / 2 (2n+I)! .
The inequality was derived by F.W. Bessel in 1828 Bessel's interpolation formula has certain advantages
for the trigonometric system. over Gauss' formulas (1), (2); in particular, if the inter-
References polation is at the middle of the segment, i.e. at t = 1 /2,
[I] KUDRYAVTSEV, L.D.: Mathematical analysis, 2, Moscow, 1973 all coefficients at the differences of odd orders vanish.
(in Russian). If the last term on the right-hand side of (3) is omitted,
L.P. Kuptsov
the polynomial B 2n + I (x 0 + th), which is not a proper
Editorial comments. Usually, the orthogonal system of
interpolation polynomial (it coincides with f (x) only in
elements {<Pa} is orthonormalized, i.e. one sets
the 2n nodes Xo -en -1)h, ... ,Xo +nh), represents a
!f;a =<Pa / II <Pa II Bessel's inequality then takes the form
better estimate of the residual term (d. Interpolation
~ I (t, !f;a) I .;;;; II f 11 2 , formula) than the interpolation polynomial of the same
aEA
degree. Thus, for instance, if x = X 0 + th EO (x 0, XI)' the
which is easier to remember. In this form i is used in
approximation theory, Fourier analysis, the theory of orthog-
estimate of the last term using the polynomial which is
onal polynomials, etc. most frequently employed
References
[A1) YOSIDA, K.: Functional analysis, Springer, 1980.
[A2) CHENEY, E.W.: Introduction to approximation theory, Chel-
sea, reprint, 1982. wri tten wi th respect to the nodes
[A3) DAVIS, P.l.: Interpolation and approximation, Dover, reprint, x0 - h, x 0, x 0 + h, x 0 + 2h, is almost 8 times better than
1975.
that of the interpolation polynomial written with
[A4) HEWITT, E. and STOMBERG, K.: Real and abstract analysis,
Springer, 1965. respect to the nodes Xo -h, xo, Xo +h or
xo,xo+h, xo+2h ([2]).
AMS 1980 Subject Classification: 42CXX, 47 A30
References
[IJ BEREZIN. IS and ZHIDKOY. N.P.: Computing methodl" L Per-
BESSEL INTERPOLATION FORMULA - A formula gamon Press. 1973 (translated from the Russian).
which is defined as half the sum of the Gauss formula [2J BAKHYALOY, N.S.: Numerical methodl': analvsis. algebra, ord,
(d. Gauss interpolation formula) for forward interpola- narr differential equations. Mir. 1977 (translated from the Rus-
sian).
tion on the nodes M.K. Samarin

374
 BEST APPROXIMATION

Editorial comments. Editorial comments. The function Kp(z) is usually called

References the modified Bessel function of the third kind.
[A 1] ABRAMOWITZ, M. and STEGUN, I.A.: Handbook of mathemati- AMS 1980 Subject Classification: 31815
cal functions, Nat. Bur. Stand., Appl. Math. Ser., 55, Dover,
1970.
[A2] HILDEBRAND, F.: Introduction to numerical analysis, BESSEL SYSTEM - A concept in the theory of
Addison-Wesley, 1956. orthogonal systems. Let {!{In} and {gn} be two complete
AMS 1980 Subject Classification: 65D05 systems of functions in L 2(a, b)=L2 (i.e. measurable
functions that are square-integrable on the segment
BESSEL POTENTIAL - A potential of the form [a, b D, forming a biorthogonal system of functions. The
system {!{In} is said to be a Bessel system if, for any
Pa(x) = jGa(x-y)dp.(y), a>O,
function fEL2' the series
R'
00
where X=(Xh'" ,xn), y=(Yh'" ,Yn) are points in ~c;
the Euclidean space Rn, n ;;'2; dp. is a Borel measure on n=1

Rn; is convergent; here, Cn=(j, gn) are the coefficients of

G a(x)=2(2-n-a)/2'/T-n/2 [r [; ]r~(n-a)/2(' x I) 1x l(a-n)/2,
the expansion 00

f "" ~cno/n

r
n=1

1 x1 = [i~ 1 xT 1 /2,
of the function fwith respect to the system {o/n}. For a
system {!{In} to be a Bessel system it is necessary and
and K p(z) is the modified cylinder function (or Bessel sufficient that it be possible to define a bounded linear
function, cf. Cylinder functions) of the second kind of operator A on the space L2 such that the system {cf>n}
order P or the Macdonald function of order P; Ga(x) is defined by the equation A !{In =cf>n (n = I, 2, ... ) is a
called a Bessel kernel. complete orthonormal system. If the system {!{In} is a
The principal properties of the Bessel kernels Ga(x) Bessel system, there exists a constant M such that for
are the same as those of the Riesz kernels (d. Riesz any fEL2 ao

potential), viz., they are positive, continuous for X=F0, ~(f, gnf ~ MIl/ilL
n=1
can be composed
References
jGa(x-y)Gp(y)dy = Ga+p(x), [1) KACZMARZ, S. and STEINHAUS, H.: Theorie der Orthogonal
R' reihen, Chelsea, reprint, 1951.
P.I. Lizorkin
but, unlike the Riesz potentials, Bessel potentials are
applicable for all 0:>0, since AMS 1980 Subject Classification: 42CXX

Ga(x) "" 2(I-n-a)/271 I- n )/2 [r (;]] -II X l(a-n-I)/2e-lxl, BEST APPROXIMATION of a function x(t) by func-
tions u(t) from a fixed set F - The quantity
as I x I~oo. E(x, F) = inf p.(x, u),
If 0:>2m, where m is a natural number, and the ueF
measure dp. is absolutely continuous with square- where p.(x, u) is the error of approximation (see
integrable density f(y)EL 2 (R2m ), the Bessel potentials Approximation of functions, measure of). The concept of
satisfy the identities: a best approximation is meaningful in an arbitrary
(l-~rPa(x) = P a - 2m (x), metric space X when p.(x, u) is defined by the distance
and between x and u; in this case E(x, F) is the distance
(l-~rp2m(X) = I(x), from x to the set F. If X is a normed linear space, then
where Ll is the Laplace operator on R2m. In other for a fixed Fe X the best approximation
words, the function G2m (x) is a fundamental solution E(x, F) = ueF
inf II x-u II (1)
of the operator (1- Ll)m .
may be regarded as a functional defined on X (the
References
[I) NIKOL'SKIi, S.M.: Approximation offunctions of several variables
functional of best approximation).
and imbedding theorems, Springer, 1975 (translated from the The functional of best approximation is continuous,
Russian). whatever the set F. If F is a subspace, the functional of
(2) ARONSZAJN, M. and SMITH, K.T.: 'Theory of Bessel potentials
1', Ann. Inst. Fourier (Grenoble) 11 (1961), 385-475.
best approximation is a semi-norm, i.e.
E.D. Solomentsev E(xI + X2, F) ~ E(Xlo F)+ E(X2, F)

375
BEST APPROXIMATION
and
E(Xx, F) = I A IE(x, F)
for any AER. If F is a finite-dimensional subspace, then
for any x EX there exists an element Uo EF (an element
of best approximation) at which the infimum in (1) is
attained:
E(x, F) = II x -Uo II
In a space X with a strictly convex norm, the element
of best approximation is unique.
Through the use of duality theorems, the best
approximation in a normed linear space X can be
expressed in terms of the supremum of the values of
certain functionals from the adjoint space X (see, e.g.
[5], l8]). If F is a closed convex subset of X, then for
any XEX
E(x,F) = J,!1f [r(x)-~~y/(u)];
II!II<;;]
in particular, if F is a subspace, then
E(x, F) = s~ I(x),
{!II<;;]
where Fl. is the set of functionals fin x' such that
f(u)=O for any UEF. In the function spaces Cor L p '
the right-hand sides of (2) and (3) take explicit forms
depending on the form of the linear functional. In a
Hilbert space H, the best approximation of an element
x EH by an n-dimensional subspace Fn is obtained by
orthogonal projection on Fn and can be calculated; one
has:
G(x, u], ... ,un)
E(x, Fn) =
G(uJ, ... ,un) ,
where U1, . . . ,Un form a basis of Fn and
G(u), ... ,un) is the Gram determinant, the elements
of which are the scalar products (Uj, Uj), i, j = 1, ... ,n.
If {Uk} is an orthonormal basis, then
E2(X, Fn) = I X 112_ ~(x, ud
k=]
In the space C = qa, b] one has the foU) wing esti-
mate for the best uniform approximation of a function
X(t)EC by an n-dimensional Chebyshev subspace
Fn cC (the de la Vallee-Poussin theorem): If for some
function U(t)EFn there exist n + 1 points tk,
a~tl < ... <tn+1 ~b, for which the difference
~(t) = x(t) - u(t)
takes values with alternating signs, then
E(x, Fn);;;' min I ~(tk) I
]<;;k<;;n+]
For best approximations in L 1(a, b) see Markov cri-
terion. In several important cases, the best approxima-
tions of functions by finite-dimensional subspaces can
be bounded from above in terms of differential-
difference characteristics (e.g. the modulus of con-
tinuity) of the approximated function or its derivatives.
376
The concept of a best uniform approximation of con-
tinuous functions by polynomials is due to P.L. Che-
byshev (1854), who developed the theoretical founrla-
tions of the concept and established a criterion for
polynomials of best approximation in the metric space
C (see Polynomial of best approximation).
The best approximation of a class of functions is the
supremum of the best approximations of the functions
f in the given class WC by a fixed set of functions F, i.e.
the quantity
E(IJH, F) = suoE(j, F) = suoinfJ.t(f, cp).
!E~ !E~4>EF
The number E(WC, F) characterizes the maximum devi-
ation (in the specific metric chosen) of the class WC
from the approximating set F and indicates the
minimal possible error to be expected when approxi-
(2)
mating an arbitrary function f E WC by functions of F.
Let WC be a subset of a normed linear function space
={
X, let U U 1(t), U 2 (t), . . .} be a linearly independent
(3) system of functions in X and let Fn, n = 1, 2, ... , be
the subspaces generated by the first n elements of this
system. By investigating the sequence E(WC, F n ),
n = 1, 2, ... , one can draw conclusions regarding both
the structural and smoothness properties of the func-
tions in WC and the approximation properties of the sys-
tem U relative to WC. If X is a Banach function space
and U is closed in X, i.e. U Fn =X, then E(WC, Fn)~O
as n ~ 00 if and only if WC is a compact subset of X.
In various important cases, e.g. when the Fn are sub-
spaces of trigonometric polynomials or periodic splines,
and the class WC is defined by conditions imposed on
the norm or on the modulus of continuity of some
derivative t'l, the numbers E(WC, Fn) can be calculated
explicitly [5]. In the non-periodic case, results are avail-
able concerning the asymptotic behaviour of E(WC, Fn)
as n~oo.
References
[1] CHEBYSHEV, P.L.: Complete collected works, 2, Moscow, 1947
(in Russian).
[2] ACHIEZER, N.!. [N.I. AKHIEZER]: Theory of approximation, F.
Ungar, 1956 (translated from the Russian).
[3] DZYADYK, V.K.: Introduction to the theory of uniform approxi-
mation offunctions by polynomials, Moscow, 1977 (in Russian).
[4] GONCHAROV, V.L.: The theory of interpolation and approxima-
tion of functions, Moscow, 1954 (in Russian).
[5] KORNEicHUK, N.P.: Extremal problems in approximation theory,
Moscow, 1976 (in Russian).
[6] NIKOL'SKIi, S.M.: Approximation offunctions of several variables
and imbedding theorems, Springer, 1975 (translated from the
Russian).
[7] TIMAN, A.F.: Theory of approximation of functions of a real
variable, Pergamon, 1963 (translated from the Russian).
[8] TIKHOMIROV, V.M.: Some questions in approximation theory,
Moscow, 1976 (in Russian).
[9] LAURENT, PJ.: Approximation et optimisation, Hermann, 1972.
N.P. Kornelchuk
VP. Motorny!

Editorial comments. In Western literature an element, a
functional or a polynomial of best approximation is also
called a best approximation.
References
[A1] LORENTZ, G.G.: Approximation offunctions, Holt, Rinehart
and Winston, 1966.
[A2] CHENEY, E.W.: Introduction to approximation theory, Chel-
sea, reprint, 1982.
[A3] RICE, J.R.: The approximation of functions, 1. Linear theory,
Addison-Wesley, 1964.
[A4] PINKUS, A: n-widths in approximation theory, Springer,
1985.
AMS 1980 Subject Classification: 41 A50, 41 A65,
41A46
BEST APPROXIMATION IN THE MEAN - The best
approximation of a function x by functions u from a
fixed set F when the measure (error) of approximation
is expressed in terms of an integral metric (see Best
approximation; Approximation in the mean).
N.P. Kornelchuk
VP. Motornyr
AMS 1980 Subject Classification: 41 A50
BEST APPROXIMATIONS, SEQUENCE OF - A
sequence {E(x, F n )}, n = 1, 2, ... , of numbers, where
E(x, Fn) is the best approximation of an element x of a
normed linear space X by elements of an n-dimensional
subspace Fn ex, with Fl eF2 e .. " so that
E(x, Fd~E(x, F2)~ .. '. Usually, Fn is the linear
span of the first n elements of some fixed system
{u 1, u2, ... } of linearly independent elements of X.
In the case X = q a, b] and Fn = F~ is the subspace
of algebraic polynomials of degree n - 1, sequences of
best approximations were first considered in the
eighteen-fifties by P.L. Chebyshev; the fact that
E(x, F~} ~ 0 as n~oo
for any function X(t)E qa, b] was established in 1885
by K. Weierstrass. In the general case, the relation
lim E(x, Fn}
n .... oo
=0 for all x EX
is always satisfied when the union of the subspaces F n ,
n = I, 2, ... , is everywhere dense in X,
UFn = X
(essentially, this is an equivalent statement). However,
the sequence {E(x, Fn)} may converge to zero arbi-
trarily slowly. This follows from a theorem of Bern-
stein: If {Fn} is a sequence of subspaces of dimension
n = 1, 2, ... , of a normed linear space X, such that
FIe F 2 e . .. and U Fn = X, then, for any monotone
decreasing null sequence {JLn} of non-negative real
numbers, there exists an XEX such that E(x, Fn)=/Ln,
n = 1, 2, .... In the function spaces C and Lp ' the rate
at which a sequence of best approximations tends to
BEST COMPLETE APPROXIMA nON
zero depends both on the system of subspaces Fn and
on the smoothness characteristics of the approximated
function x (the modulus of continuity, the existence of
derivatives up to a specific order, etc.), and it can be
estimated in terms of these characteristics. Conversely,
knowing the rate of convergence to zero of the
sequence {E(x, Fn )}, one can draw conclusions with
respect to the smoothness of ,x (t) (see Approximation of
functions, direct and inverse theorems).
References
[1] BERNSTEIN, S.N. [S.N. BERNSHTEiN]: Collected works, 2, Mos-
cow, 1954 (in Russian).
[2] GoNCHAROV, V.L.: The theory of interpolation and approxima
tion offunctions, Moscow, 1954 (in Russian).
[3] TIMAN, AF.: Theory of approximation offunctions of a real
variable, Pergamon, 1963 (translated from the Russian).
N.P. Kornerchuk
VP. Motornyr
Editorial comments. Theorems inferring smoothness
characteristics of a function x E C or Lp from properties of
E(x, Fn) were first given by D. Jackson in 1911 for algebraic
or trigonometric approximation, see Jackson theorem.
Theorems converse to these, i.e. inferring properties of
E(x, Fn} from smoothness characteristics of the function x,
have been proved by S.N. Bernstein [S.N. BernshteTn] and
A. Zygmund, cf. Bernstein theorem. See also [A2], Chapt.
4, Sect. 6 and Chapt. 6, Sect. 3.
References
[A1] NATANsoN, J.P.: Constructive function theory, 1-3, F. Ungar,
1964-1965 (translated from the Russian).
[A2] CHENEY, E.W.: Introduction to approximation theory, Chel-
sea, reprint, 1982.
AMS 1980 Subject Classification: 41 A50
BEST COMPLETE APPROXIMATION - A best
approximation of a function f(Xl, ... ,Xk) in k vari-
ables (k~2) by algebraic or trigonometric polynomials.
Let X be the space C or ~ of functions f(XIo ... ,Xk),
217-periodic in each variable, that are either continuous
or p-summable (p ~ 1) on the k-dimensional period
cube with edges of length 217.
The best complete approximation of a function
f(Xh ... ,Xk)EX by trigonometric polynomials is the
quantity
En, . ... .n, (f)x = T,,~f,. II f - Tn, .. .. .n, II x,
where the infimum is taken over all trigonometric poly-
nomials of degree ni in the variable Xi (1 ~ i ~k).
Together with the best complete approximation, one
also considers best partial approximations.
A best partial approximation of a function
f(XIo . .. ,XdEX is a best approximation by functions
Tn." ... ,n., (x 1, . . . ,Xk) E K that are trigonometric poly-
nomials of degree n v" .. ,n v, (l~r<k), respectively,
in the fixed variables x v,, ... ,x v, with coefficients
377
BEST COMPLETE APPROXIMATION
depending on the remaining k - r variables, i.e.
En,!' ... ,n",oo(f)x =~ inf IIf-Tn,!, ... ,n"llx.
11'1 .. ..,.
It is obvious that
S.N. Bernstein [S.N. Bernshtein] [1] proved the fol-
lowing inequality for continuous functions in two vari-
ables:
(1)
.;;; A In(2 + min { n J, n2 })(En"oo(f)c + En"oo(f)c),
where A is an absolute constant. It has been shown [3]
that the term In(2+min{nbn2}) in inequality (1) (and
in the analogous relation for the space L I) cannot be
replaced by a factor with a slower rate of increase as
min{nb n2}-HX).
In the space Lp (p> 1) one has the inequality
k l~';;;;oo, of 2'7T-periodic functions); let F=Tn be the
En" ... ,n,(f)Lp .;;; Ap,k ~Eni,oo(f)Lp' (2)
i=l
where the constant Ap,k depends only on p and k.
Similar definitions yield best complete approxima-
tions and best partial approximations of functions
defined on a closed bounded domain n C Rk by alge-
braic polynomials, and in this case inequalities similar
to (1) and (2) have been established.
References
[1) BERNSHTEiN, S.N.: Collected works, 2, Moscow, 1954 (in Rus-
sian).
(2) TIMAN, A.F.: Theory of approximation offunctions of a real
variable, Pergamon, 1963 (translated from the Russian).
(3) TEMLYAKov, V.N.: 'On best approximations of functions in
two variables', Dokl. Akad. Nauk SSSR 223, no. 5 (1975),
1079-1082 (in Russian). NP v ~ h k
. . .n.ornelC U
v.P. Motornyr
Editorial comments.
References
[A 1) LORENTZ, G.G.: Approximation of functions, Holt, Rinehart
and Winston, 1966. ) same error (over the entire class) in the metric of C
AMS 1980 Subject Classification: 41 A50, 41 A63
BEST LINEAR METHOD - With respect to the
approximation of elements in a given set m, the linear
method that yields the smallest error among all linear
methods. In a normed linear space X, a linear method
for the approximation of elements x Em C X by ele-
ments of a fixed subspace Fe X is represented by a
linear operator that maps the entire space X, or some
linear manifold containing m, into F. If !> is the set of
all such operators, a best linear_ method for m (if it
exists) is defined by an operator A E!> for which
.~~~ I x - Ax II = }~t .~~~ I x - Ax II
The method defined by an operator A in !> will cer-
tainly be a best linear method for m relative to the
378
approximating set F if, for all x Em,
II x-Ax I .;;; ~~gE(x, F)
(E(x, F) is the best approximation of x by F) and if,
moreover, for all x EX,
I x-Ax II = E(x, F).
The latter is certainly true if X is a Hilbert space,
F = Fn IS an n-dimensional subspace of X,
n = 1, 2, ... , and A is the orthogonal projection onto
Ax = ~ (x, ek)eb
k=l
where {e I, . . . ,en} is an orthonormal basis in Fn.
Let X be a Banach space of functions defined on the
entire real line, with a translation-invariant norm:
II X('+T) II = II xO II (this condition holds, e.g. for the
norms of the spaces C=C[O,2'7T] and Lp =Lp(0,2'7T),
subspace of trigonometric polynomials of order n.
There exist best linear methods (relative to Tn) for a
class m of functions X(t)EX that contains x(t+a) for
any aER whenever it contains x(t). An example is the
linear method
/loao
A(x;t;/L,p) = -2-+ (*)
n
+ ~ (/Lk(ak cos kt + bk sinkt)+ Pk(ak sin kt - bk coskt)},
k=l
where ak and bk are the Fourier coefficients of x(t)
relative to the trigonometric system, and ILk and Pk are
numbers.
Now consider the classes W'aoM (and WiM),
r = 1, 2, ... , of 2'7T-periodic functions x(t) whose
derivatives x(r-l)(t) are locally absolutely continuous
and whose derivatives x(r)(t) are bounded in norm in
Lao (respectively, in L I) by a number M. For these
classes, best linear methods of the type (*) yield the
(respectively, Ld as the best approximation by a
subspace Tn; the analogous assertion is true for these
classes with any rational number r>O (interpreting the
derivatives x(r)(t) in the sense of Weyl). For integers
r = 1, 2, ... , best linear methods of type (*) have been
constructed using only the coeff~cients ILk (all Pk = 0).
If F = S~ is the subspace of 2'7T-periodic polynomial
splines of order r and defect 1 with respect to the parti-
tion kw / n, k = 0, +- 1, ... , then a best linear method
for the classes W'ao+ I M (and Wi + I), r = 1, 2, ... , is
achieved in Lp ' l:!(p';;; 00 (resp. in L I) by splines in S~
interpolating the function x(t) at the points
kw /n+[l+(~IY1'7T /4n .
References
[1) ACHIEZER, N.I. [N.I. AKHIEZER): Theory of approximation, F.
Ungar, 1956 (translated from the Russian).
 BETA-DISTRIBUTION

[2] KORNEicHUK, N.P.: Extremal problems in approximation theory,
Moscow, 1976 (in Russian).
[3] TIKHOMIROV, Y.M.: Some questions in approximation theory,
Moscow, 1976 (in Russian).
N.P. Korneichuk
v.P. Motornyr
Editorial comments.
References
[A 1] KrnSEWETIER, H.: Vorlesungen Dber lineare Approximation,
Deutsch. Verlag Wissenschaft., 1973.
[A2] RICE, J.R.: The approximation of functions, 1. Linear theory,
Addison-Wesley, 1964.
AMS 1980 Subject Classification: 41 A65, 41 A50
BEST QUADRATURE FORMULA, optimal quadra-
ture formula - An approximate integration formula that
guarantees the minimum error for a given class of func-
tions, relative to all formulas of a specified type. As an
example, consider the quadrature formula
b n m A few recent textbooks are listed below.
jp(x)f(x)dx = ~ ~Pk,f(I)(Xk)+R(f), (*)
a k=li=O
where p(x) is a weight function. The remainder (error)
term R(j)=R(j, Xn, Pnm ) depends both on the func-
tion f(x), and on the vector (Xn' P nm ) consisting of the
interpolation nodes Xk (it is usually assumed that
Xk E[a, b D and the coefficients hi, k = 1, ... ,n;
i = 0, ... , m. Fixing n ~ I and m ~O, let A denote
some set of vectors (Xn, P nm ) (and hence also some set
of quadrature formulas), defined by some restrictions
on the interpolation nodes and coeffi~ents (in particu-
lar, one might consider the ~et A =A(Xn) of coefficients
hi for a fixed node vector Xn). Let we be some class of
functions f(x), it being assumed that the integral and
the sum in (*) exist. The best quadrature formula of
type (*) for the class WC relative to the set A is defined
drature formula generalizes in a natural way to func-
tions of several variables (cubature formulas).
References
[I] NIKOL'SKli,S.M.: Quadratureformulae, H.M. Stationery Office,
London, 1966 (translated from the Russian).
[2] KRYLOV, N.M.: Approximate calculation of integrals, Macmil-
lan, 1962 (translated from the Russian).
[3] LAURENT, PJ.: Approximation et optimisation, Hermann, 1972.
[4] ZHENSYKBAEV, A.A.: 'Monosplines of minimal norm and qua-
drature formulas', Uspekhi Mat. Nauk 36, no. 4 (1981), 107-159
(in Russian). N P v "h k
. . n.ornelC U
v.P. Motornyr
Editorial comments. The terminology 'best formula' is
often encountered in the literature on numerical analysis,
but, as was obseNed in [A2J, p. 75, it should be taken with
a large dose of salt, because, after all, any quadrature for-
mula, no matter how the weights Pkl and the nodes Xk are
chosen, will exactly integrate an infinite-dimensional family
of functions.
References
[A 1] BRASS, H.: Quadraturverfahren, Vanden hoek & Ruprecht,
1977.
[A2] DAVIS, P.J. and RABINOWITZ, P.: Methods of numerical
integration, Acad. Press, 1984.
[A3] ENGELS, H.: Numerical quadrature and cubature, Acad.
Press, 1980.
AMS 1980 Subject Classification: 65030, 65032
BETA-DISTRIBUTION - A continuous probability
distribution concentrated on (0, 1) with density
f3 m,n (x) = 1
B(rn, n)
xm-I(I- xt- I
,
(I)
where the parameters m, n are non-negative and the
normalizing factor B (m, n) is Euler's beta-function
1

by a vector (x:, P~m) for which

B(
rn, n
) = jx m-I(I_ x )n-I dx = f(rn)f(n)
0 rem + n) ,
sup I R(f, X~, P~m)
f c '1fl
I= where r(n) is the gamma-function. The distribution
function is expressed as the incomplete beta-function

The constmction of best quadrature formulas is mtI- Bm.n(x) = 1 j'ym-I(I_y)"-1 dy, O<x<1
B(rn, n) 0
mately connected with certain problems in spline
approximation; in many cases it reduces to minimizing (this function has been tabulated, see [1], [2]). The
the norm of a monospline (see [ID, Best quadrature for- moments of the beta-distribution are given by the for-
mulas
mulas, together with sharp estimates for the remainder _ B(rn+k, n) _
rnk - B(rn, n) , k-I,2, ....
term, are known for many important classes of continu-
ous and differentiable functions. From a more general In particular, the mathematical expectation and the
point of view, the problem of finding best quadrature variance are m / (m + n) and
formulas and the corresponding errors for a class we mn / {(m+n)2(m+n+ I)}, respectively. If m>l and
may be viewed as the problem of optimal recovery of a n > L the density curve f3m.n(x) has a single mode at
functional b the point x =(m -1) / (m +n -2) and vanishes at the
1(f) = jp(x)f(x)dx, ends of the interval. If either m < 1 or n < 1, one ordi-
nate at the end of the graph becomes infinite, and if
where fEWC, on the basis of the information {fU)(xd} , both m < I and n < 1, both ordinates at the ends of the
k = 1, ... ,n; i =0, ... ,m. The concept of a best qua- interval are infinite and the curve is U-shaped. If m = 1

379
BET A-DISTRIBUTION

and n = I the beta-distribution reduces to the uniform abIes p and q which, for. p, q >0, is defined by the
distribution on the interval (0, I). Another special case equation I
of the beta-distribution is the so-called arcsine distribu- B(p, q) = jxP-I(l-X)q-1 dx. (*)
tion: o
I
f31/2,1/2(X) = ?TVx(l-x) The values of the beta-function for various values of
the parameters p and q are connected by the following
If one substitutes x = I / (I + t) in (I), then one obtains
relationships:
a distribution having the density B(P, q) = B(q,p),
1 1m - q-I
(J'
fJ (I) -
I
0<1<00, (2) B(P,q) =
p
+ q- \B(p,q-I), q>l.
m,n - B(m, n) (1 +t)m+n-2 '

This distribution is called a beta-distribution of the The following formula is valid:

second kind, as distinct from the beta-distribution (l). ?T
B(p, I-p) = -.-,
smp?T
O<p<l.
The distributions (l) and (2) correspond to 'type I' and
'type VI' distributions in the system of Pearson curves, If P and q are complex, the integral (*) converges if
An important case of generation of a beta-distribution Rep>O and Req>O. The beta-function can be
is the following: If X I and X 2 are independent and expressed by the gamma-function:
have gamma-distributions (cf. Gamma-distribution) with B(P q) = f(P)f(q)
respective parameters m and n, then the random vari- , f(P +q)
able X I / (X I + X 2) will have a beta-distribution with VI. Bityutskov
density f3m,n(x). This fact to a large extent explains the AMS 1980 Subject Classification: 33A 15
role played by beta-distributions in various applica-
tions, in particular in mathematical statistics: The dis- BETTI GROUP - In a broad sense, the same as a
tributions of several important statistics are reducible homology group; in a narrow sense, the Betti group is
to beta-distributions. For instance, the distribution the free part of the homology group with as domain of
function of the F-relationship coefficients the group Z of integers, if this homology
group is finitely generated. Named after E. Betti (1823
- 1892).
References
(the random variable X~ has a i-distribution with k [IJ SEIFERT. H. and THRELFALL. W.: Lehrbuch der Topologie. Chel-
degrees of freedom) is expressed by the formula sea, reprint, 1980.
[21 ALEKSANDROV. P.S.: An introduction to homological dimension
theory and general combinatorial topology. Moscow. 1975 (in
'
P(Fmn<X)=Bm/2n/2r
' l n +mxmx J Russian).
M. 1. VoUsekhovskil
(the values of the F-distribution are usually calculated Editorial comments.
with the aid of tables of beta-functions). The beta- References
distribution function also allows one to compute the [A 11 SPANIER. E.H.: Algebraic topology, McGraw-Hili, 1966.
values of the binomial distribution functions, in view of AMS 1980 Subject Classification: 55NXX
the relationship
BETTI NUMBER, r-dimensional Betti number pr of a
= k~O [z )pk(1_p)"-k
m )
Bn-m,nd l(l-P) complex K - The rank of the r-dimensional Betti group
with integral coefficients. For each r the Betti number
Beta-distributions are used m fields other than pr is a topological invariant of the polyhedron which
mathematical statistics; thus, the density of the beta- realizes the complex K, and it indicates the number of
distribution is the weight function for the system of pairwise non-homological (over the rational numbers)
orthogonal Jacobi polynomials. cycles in it. For instance, for the sphere sn:
References pO = I. pi = ... =p,,-I = O. p" = I;
[lJ BOL'SHEV, L.N. and SMIRNOV. N,V.: Tables of mathematical
statistics, Libr. of mathematical tables, 46. Nauka, Moscow, for the projective plane P2(R):
1968 (in Russian). Processed by L.S. Bark and E.S. Kedova.
[2J PEARSON, K.: Tables of the incomplete beta:function, Cambridge po = I. pi = p2 = 0;
U niv, Press, 1932.
A. V Prokhorol' for the torus T2:
AMS 1980 Subject Classification: 62EXX p" = p2 = I. pi = 2.
For an n-dimensional complex K" the sum
BETA-FUNCTION, Bjunction. Euler Bjunction.
Euler integral of the first kind - A function of two vari- ~(-llP'
, 1I

380

is equal to its Euler characteristic. Betti numbers were
introduced by E. Betti [1].
References
(1) BETII, E.: Ann. Mat. Pura Appl. 4 (1871), 140-158.
M.1. Voitsekhovskii
Editorial comments.
References
[A1] SPANIER, E.H.: Algebraic topology, McGraw-Hili, 1966.
AMS 1980 Subject Classification: 55NXX
BEZOUT RING - An integral domain with a unit ele-
ment in which any ideal of finite type is principal. Any
principal ideal ring and also any valuation ring is
Bezout. A Bezout ring is integrally closed, and its local-
izations (i.e. its rings of fractions with respect to multi-
plicative systems S, cf. Localization in commutative
algebra) are again Bezout rings. For a finite set
ai, ... ,an of elements of a Bezout ring A there exist
the greatest common divisor (the greatest common divi-
sor of (ai, ... ,an) has the form ~bA, b;EA, a so-
called Bezout identity) and the lowest common multiple.
A Noetherian ring (and even a ring that satisfies the
ascending chain condition only for principal ideals)
which is Bezout is a principal ideal ring. As for princi-
pal ideal rings, a module of finite type over a Bezout
ring is a direct sum of a torsion module and a free
module.
VI. Danilov
Editorial comments.
References
[A1] GILMER, R. and DEKKER, M.: Multiplicative ideal theory,
1972.
AMS 1980 Subject Classification: 13G05, 13F99,
13F10
BEZOUT THEOREM - 1) Bezout's theorem on the
division of a polynomial by a linear binomial: The
remainder of the division of the polynomial
j(x) = aox'+ ... +a.
by the binomial x - a is f (a). It
is assumed that the
coefficients of the polynomials are contained in a cer-
tain commutative ring with a unit element, e.g. in the
field of real or complex numbers. A consequence of
Bezout's theorem is the following: A number a is a root
of the polynomial f (x) if and only if f (x) is divisible
by the binomial x -a without remainder.
2) Bezout's theorem on homogeneous equations: If a
system of n homogeneous equations in n + 1 unknowns
j;(xo, ... ,x.) = 0, ;=1, ... ,n, (*)
has only a finite number of non-proportional non-zero
solutions in an algebraically closed field containing the
coefficients of the system, then the number of these
solutions counted according to their multiplicity is
BIANCHI IDENTITY
equal to the product of the degrees of the equations.
The multiplicity of the solutions is, by definition, the
intersection index (in algebraic geometry) of the hyper-
surfaces (*) at the respective point. The theorem is
called after E. Bezout [1], who studied systems of alge-
braic equations of higher degrees.
References
[1) BEZOUT, E.: Theorie generale lies equations algebriques, Paris,
1779.
VN. Remeslennikov
V.E. Voskresenskii
Editorial comments.
References
[A1] SHAFAREVICH, I.R.: Basic algebraic geometry, Springer,
1977, Chapt. 4, Sect. 2 (translated from the RUSSian).
AMS 1980 Subject Classification: 14C17, 12E05
BIANCHI CONGRUENCE, B-congruence A
congruence of straight lines in which the curvatures of
the focal surfaces at the points situated on the same
straight line of the congruence are equal and negative.
The principal surfaces of a B-congruence cut out conju-
gate line systems on its focal surfaces. The straight
lines of the congruence map the asymptotic nets of the
focal surfaces onto an orthogonal net on a sphere. The
curvature of a focal surface of a Bianchi congruence is
expressed in asymptotic parameters u and v by the for-
mula:
K= 1
(</>(u)+tf;(V2 .
Surfaces whose curvatures satisfy this condition are
called Bianchi surfaces (B-surfaces, cf. Bianchi surface).
Bianchi congruences were studied by L. Bianchi [1].
References
[I) BIANCHI, L.: Ann. Mat. Pura Appl. 18 (1890), 301-358.
[2] FINIKOV, S.P.: Theorie der Kongruenzen, Akademie-Verlag,
1959 (translated from the Russian).
[3) FINIKOV, S.P.: Bending and related geometrical problems,
Moscow-Leningrad, 1937 (in Russian).
[4] ZHULIKOVSKIi, V.I.: Classical differential geometry in a tensor
setting, Moscow, 1963 (in Russian).
VT. Bazylev
AMS 1980 Subject Classification: 53A25, 53A05
BIANCHI IDENTITY - A relation connecting the
components of the covariant derivatives of the curva-
ture tensor R;Jk of a Riemannian space:
Ri'k.l + Ri%,.j + Ri~j.k = 0,
where h, i, j, k, 1= 1, ... , n. First established by L.
Bianchi [1] in 1902.
References
[I] BIANCHI, L.: Lezioni di geometria differenziale, 1-2, Zanichelli,
Bologna, 1923-1927.
M.l. Voftsekhovskif
Editorial comments. Here Ri7k,l denotes of course the
covariant derivative of ROk with respect to the /-th coordi-
nate.
381
BIANCHI IDENTITY

The identity described above is often called the second Bianchi congruence). If S is a pseudo-sphere then S'
Bianchi identity. The first Bianchi identity is then given by also is a pseudo-sphere. The pseudo-spheres S' which
Rjk,+Rk'J+R/ik =0 are Bianchi transforms of S are orthogonal trajectories
(see [A1], [A2]). Generalized versions of these identities for of the following congruence of circles. They are
curvature forms and curvature tensors of connections with situated in the tangent plane to S and have the same
torsion are given in [A2]. radius as S.
VT. Bazylev
References Editorial comments. Cf. also [A2], articles 803, 804 in
[A1] HICKS, N.J.: Notes on differential geometry, v. Nostrand,
1965. volume III.
[A2] KOBAYASHI, S. and NOMIZU, K.: Foundations of differential References
geometry, 1, Wiley (Interscience), 1963. [A 1]EISENHART, L.P.: A treatise on the differential geometry of
curves and surfaces, Boston, 1909.
AMS 1980 Subject Classification: 53820
[A2] DARBOUX, G.: Ler;;ons sur la theorie generale des surfaces,
Chelsea, reprint, 1972.
BIANCHI SURFACE - A surface of negative Gauss-
AMS 1980 Subject Classification: 53A05, 53A25
ian curvature K which can be expressed in asymptotic
parameters (u, v) as BIASED ESTIMATOR - A statistical estimator
K= whose expectation does not coincide with the value
[U(u)+ v(v)f ' being estimated.
where U(u) and V(v) are arbitrary functions; thus, Let X be a random variable taking values in a sam-
Bianchi surfaces can be characterized by the fact that pling space (1:, !Jd, PII ), OEe, and let T= T(X) be a sta-
the function (- K) - I /2 is diagonal with respect to its tistical point estimator of a function f(O) defined on the
asymptotic net, i.e. parameter set e. It is assumed that the mathematical
a2(-K)I/2 _ expectation Ell {T} of T exists. If the function
auav - o. b(8) = Ee{T}-f(8) = Ee{T-f(8)}
For instance, a ruled Bianchi surface is a conoid - a is not identically equal to zero, that is, b( 0) = 0, then T
surface attached to the Peterson surface. If a Bianchi is called a biased estimator of f(O) and b(O) is called the
surface is given, it is possible to determine the classes bias or systematic error of T.
of surfaces which can be obtained by a deformation Example. Let XI. ... ,Xn be mutually-independent
over a principal base and to classify them. Thus, if the random variables with the same normal distribution
principal base contains two families of geodesic lines, N I (a, 0 2 ), and let
the functions U and V are constant and the bent sur-
XI + ... +X"
face is a Voss surface (the class B 0)' X= - - - - - -
n
The class B I is characterized by the fact that only
one family of lines of the principal base are geodesics Then the statistic
(one of the functions U, V is constant); conoids may S~ = 1- (X, - X)2
serve as an example. The class B 2 corresponds to func- n,=1
tions U, V which both depend non-trivia )'y on their is a biased estimator of the variance 02 since
arguments. See also Bianchi congruence.
n-i a2
J.Kh. Sabitov E{S~} = - - a2
n
= a"'--,
1

n
Editorial comments. The notion of a 'deformation over a that is, the estimator S~ has bias b( 0 2 ) = - 0 2 / n. The
principal base' is not very common in Western literature, mean-square error of this biased estimator is
and this type of deformation has no standard name even. It
can best be characterized as a deformation preserving a E{(S~-~i} = 2n~i a4 .
n-
conjugate net.
The best unbiased estimator of 02 is the statistic
References
[A1] L Leziom di geometria differenziale. 2. Zanichelli.
BIANCHI.
Bologna, 1927. Chap!. 1.
S2 = _n_S2 = _i_(x,_X)2,
" n-i n n-i,cl
[A2] FINIKOV. S.P.: Theorie der Kongruenzen Akademle-Verlag.
1959 (translated from the Russian). with mean-square error
AMS 1980 Subject Classification. 53A05 D{s~} = E{(s;-~f} = _2_a4
n-i
BIANCHI TRANSFORMATION - The
transition When n > 2, the mean-square error of the biased esti-
from one focal surface S of a Bianchi congruence to the mator S~ is less than that of the best unbiased estima-
other focal surface S' of the same congruence (cL tor s~.

382

There are situations when unbiased estimators do not
exist. For example, there is no unbiased estimator for
the absolute value 1 a 1 of the mathematical expecta-
tion a of the normal law N I (a, a2), that is, it is only
possible to construct biased estimators for 1 a I.
References
[1] CRAMER, H.: Mathematical methods of statistics, Princeton
Univ. Press, 1946.
M.S. Nikulin
AMS 1980 Subject Classification: 62F11,
62H12
BICATEGORY - A category ~ in which some sub-
categories of epimorphisms @ and of monomorphisms
me have been distinguished such that the following con-
ditions are met:
1) all morphisms a in st are decomposable into a
product a=vJ-t, where VE@, J-tE9n;
2) if VJ-t=pT, where v, PE@, J-t, TEme, then there exists
an isomorphism () such that p= vO, and T=(}-I J-t;
n
3) @ 9n coincides with the class of isomorphisms
in the category ~.
The epimorphisms in @ (the monomorphisms in me)
are called the permissible epimorphisms (monomor-
phisms) of the bicategory.
The concept of a bicategory axiomatizes the possibil-
ity of a decomposition of an arbitrary mapping into a
product of a surjective and an injective mapping. The
category of sets, the category of sets with a marked
point and the category of groups are bicategories with a
unique bicategorical structure. In the category of all
topological spaces and in the category of all associative
rings there are entire classes of different bicategorical
structures.
References
[1] TSALENKO, M.SH. and SHUL'GEiFER, E.G.: Fundamentals of
category theory, Moscow, 1974 (in Russian). M Sh T. I
. . sa en
Editorial comments. In the literature there is much con-
fusion about the terms bicategory, double category and 2-
category. As the terminology is not settled yet the reader
should take care. In this Encyclopaedia the term bicategory
is always used as defined above.
AMS 1980 Subject Classification: 18005
BICHARACTERISTIC, bicharacteristic strip, of a
linear partial differential operator - A line of tangency
of any two characteristics (cf. Characteristic)
q,(xJ, ... ,xn) = 0, I/;{xJ, ... ,xn) =0
of this linear partial differential operator. If the param-
eter s is introduced on the bicharacteristic strip, then its
equations Xj=Xj(s), i=I, ... ,n, are defined by solving
a system of 2n ordinary differential equations
x;(S) = Q~, g; = -Qx" j= 1, ... ,n.
BICOMPLEX
Here Q(~}, ... '~n' X}, . . ,xn ) is the principal symbol
of the linear partial differential operator, the dot indi-
cates differentiation with respect to the parameter s
and, if ~j =f/>x;, the equation Q =0 is the characteristic
equation of the differential operator. Thus, the solution
Xj=Xj(s), ~j=~j(s), i=I, ... ,n, of the system (*) for
Q =0 defines the bicharacteristic strip Q =O. This
bicharacteristic strip belongs to the characteristic
62F12, q,(x}, ... ,xn)=O, i.e. q,(XI(S), ... ,xn(s=O, if the
equations
and
~;(S) = c/lx;(XI(S), ... ,xn(s, j=I, ... ,n,
are valid for at least one value of s, then it follows that
they are valid for all values of s.
References
[I] COURANT, R.: Methods of mathematical physics. Partial dif-
ferential equations, 2, Interscience, 1965 (translated from the
German). B.L Rozhdestvenskit
Editorial
=
comments. The projections x;=x;(s),
i 1, ... , n, into x-space are called the bicharacteristic
curves (or rays). The bicharacteristic curves are tangent to
the characteristic hypersurfaces q,(X1, ... ,xn)=O due to
the homogeneity of the principal symbol, as a function of
(~1' ... ,~n), of degree equal to the order of the linear par-
tial differential operator (cf. also Principal part of a dif-
ferential operator; Symbol of an operator).
Nowadays, the standard reference on these matters is
[A1J, or the older, more concise, [A2].
References
[A1] HORMANDER, L.: The analysis of partial differential operators,
1, Springer, 1983, 271; 302.
[A2] HORMANDER, L.: Linear partial differential operators,
Springer, 1963, 29; 31.
AMS 1980 Subject Classification: 58FXX, 58GXX
BICOMPLEX, binary complex, double complex - A
k graded module, i.e. one representable as the direct sum
0
~ m,n of its submodules A m,n, together with a pair of
differentials d I : A m,n ~ Am + I.n ,
d 2 :A m,n ~Am,n+l,
which satisfy the conditions
dld l = 0, d 2 d2 = 0, d 2 d l +d l d 2 = O.
Instead of the direct sum, the set {A m,n} and the dif-
ferentials
d I : A m,n ~ Am -I.n ,
d 2 :A m.n ~ Am.n-I,
satisfying the above conditions, may be considered.
V. E. Govorov
Editorial comments.
References
[A1] MAc LANE, S.: Homology, Springer, 1963.
(*) AMS 1980 Subject Classification: 18G35
383
BICONNECTED SPACE

BICONNECfED SPACE - A space that cannot be where O~(J<'1T, - 00<7<00. The coordinate surfaces
decomposed into a sum of two connected non- are: the family of pairs of circular cylinders with paral-
intersecting proper subsets containing more than one lel axes (7=const), the family of circular cylinders
point. orthogonal to the former (a=const), and the planes
A.A. Mal'tsev
(z =const). The system of bicylindrical coordinates is
AMS 1980 Subject Classification: 54D05 obtained as the result of translation of the system of
bipolar coordinates in the xy-plane parallel to the z-
BICYCLIC SEMI-GROUP - The semi-group with aXIS.
unit element and two generators a, b subject to the sin- The Lame coefficients (d. Lame coefficients) are:
gle generating relation ab = 1. One of the realizations of a2
the bicyclic semi-group is the Cartesian square N X N, L. = LT = (coshT-cosaf
, L.
"
= 1.
where N is the set of non-negative integers, with respect
The Laplace operator is:
to the operation
(k, 1)* (m, n) = (k + m - min(/, m), 1+ n - min(/, m)). Do! = 1
-(coshT-cosaf
a2
[naa n] naz
2
+
a~
+ 2
.
The bicyclic semi-group is an inversion semi-group and
D.D. Sokolov
as such is monogenic, i.e. is generated by a single ele-
AMS 1980 Subject Classification: 51 N20
ment. The idempotents (d. Idempotent) of the bicyclic
semi-group form a chain, which is ordered with respect
BICYLINDRICAL DOMAIN - A domain D in the
to the type of the positive numbers. The bicyclic semi-
complex space C2 that can be represented in the form
group is bisimple (cf. Simple semi-group).
of the Cartesian product of two planar domains D I
The bicyclic semi-group often occurs in theoretical
and D 2 , i.e.
investigations concerning semi-groups, not only as a
representative of certain important classes of semi- D = {(ZhZ2): zIED], Z2ED2}.
groups, but also as a 'block', which defines the struc- A special case of a bicylindrical domain is the bidisc
ture of individual semi-groups. Thus, for any idempo- (bicylinder) B(a, r)={(zj,Z2): I zl-al I <rl,
tent e of a O-simple, but not completely O-simple semi- I Z2 -a2 I <r2} of radius r=(r\, r2) with centre at
group S there exists a bicyclic sub-semi-group in S con- a=(aj,a2)' The Cartesian product of n (for n~3)
taining e as the unit element (d. [1], Para. 2.7). The ele- planar domains is said to be a polycylindrical domain. A
ments a and b of the bicyclic semi-group B defined as polydisc (polycylinder) is defined in a similar way.
above, are, respectively, its left and right multiplying M. Shirinbekov
elements (i.e. there exist proper subsets X and Y in B AMS 1980 Subject Classification: 32A07
such that aX=B, Yb =B). Moreover, in a semi-group S
with unit element e the element c will be a left multi- BICYLINDRICS - Two closed curves obtained as the
plier if and only if S contains the bicyclic semi-group intersection of two cylinders the axes of which intersect
whose unit element is identical with c; a similar at right angles. The parametric equations of bicylindrics
theorem is also valid for right multiplying elements, so are
that S has left multiplying elements if an c') only if it x = a cost. y = Vb'--a-sm-t.
,., z .
= aSlnt, b :?a,
also has right multiplying clements.
where a and b are the radii of the cylinders and t is a
References
[I] CLIFFORD, A.H. and PRESTON. G.B.: The algebraic theory of
parameter. If a = b, the bicylindrics is a pair of
semigroups, 1-2. Amer. Math. Soc .. 1961-1967. congruent ellipses.
E. V. Shikin
[2] LYAPIN. E.S.: Semigroups. Amer. Math. Soc .. 1974 (translated
from the Russian). AMS 1980 Subject Classification: 51 N20
L.N. Shevrin
AMS 1980 Subject Classification: 20MXX BIEBERBACH CONJEcruRE - A hypothesis enun-
ciated in 1916 by L. Bieberbach [I]: For all functions
BICYLINDER - A synonym for bicylindrical domain. j (z) of class S, i.e. functions j (z) which are regular
AMS 1980 Subject Classification: 32A07 and univalent in the disc I z I < I and which have the
expansIOn %

BICYLINDRICAL COORDINATES - The numbers !(z) = z + ~cnzn

2
a and z related to the rectangular Cartesian coordi-
11
T.

nates x, y and z by the formulas III this disc, one has the estimate I Cil I ~11, 11 ~2. and
I Cil I = 11 only for the Koehe jUl1c[iol1s
. _ a sinh T , _ a Sin a
~\ - COShT-cosa' ) - COShT-cosa' ~ h(z) = ;:(I-ed'z) "

3R4

where () is a real number. Bieberbach proved his con-
jecture for n = 2. The problem of finding an accurate
estimate of the coefficients for the class S is a special
case of the coefficient problem.
Owing to its simple formulation and profound signi-
ficance, Bieberbach's conjecture attracted the attention
of numerous mathematicians and stimulated the
development of different methods in the geometric
theory of functions of a complex variable. At the time
of writing (1977) the validity of the Bieberbach conjec-
ture had been established for n ";;;6. It was first proved
for n = 3 in 1923 by K. Loewner who introduced the
parametric method (cf. Parametric representation
method); other proofs of the estimate 1C3 1.,;;;3 subse-
quently appeared, and were based on variational
methods, on parametric methods and on the method of
the extremal metric. The Bieberbach conjecture for
n = 4 was first proved in 1955 by the simultaneous use
of variational and parametric methods. In 1960 the esti-
mate C4 .,;;;4 was obtained much more simply, with
1 1
the aid of Grunsky's univalence condition. This esti-
mate was also obtained by variational methods, and
with the aid of geometric reasoning; it was once more
obtained by using the Grunsky inequalities in matrix
form. The validity of Bieberbach's conjecture for n = 6
was demonstrated in 1968 with the aid of the Grunsky
inequalities; it was proved for n = 5 in 1972 by varia-
tional methods. Among other results obtained in trying
to prove the Bieberbach conjecture, the following are
worthy of note.
W. Hayman [4] obtained a number of results on the
asymptotic behaviour of the coefficients of functions
which are p-sheeted in 1 z 1 < 1, in particular for the
class S. He proved that the limit
. 1 Cn 1
hm--=a/
n-'X; n Editorial comments.
exists and that af";;; 1 with an equality sign only for the
Koebe functions.
A number of studies are available on the local
Bieberbach conjecture, i.e. the conjecture that the Koebe
function gives max Cn at least for those functions in
1 1
the class S which are close to it in the relevant
topology (cf. Univalent function). It was found that for
any n = 3, 4, ... , there exists a sufficiently small n >0
such that for a function 1 (z) E S which satisfies the con-
dition C2 - 2 <n the estimate Re cn.,;;;n is valid, and
1 1
Recn =n only for 10(z).
An estimate of the coefficients for all n, accurate as
regards the order of dependence on n, was first
obtained in 1925 by J.E. Littlewood by reducing the
estimate for the coefficients to an estimate for the
(average) integral of the modulus. More accurate esti-
mates were obtained in 1951 by I.E. Bazilevich
BIEBERBACH CONJECTURE
(I cn 1 e/2)n+const), and in 1965 by I.M. Milin
(ICn < 1.243 n, n~2).
1
The best estimates to date (1977) were obtained in
1972 [7]:
1 Cn 1 < -vf" n < 1.081 n, n ;;'2,
and in 1976 [10]:
1 Cn 1 < 1.0691 n, n;;.2.
For reviews of studies on this subject see [2], [3], [9].
References
[I] BIEBERBACH. L.: 'Uber die Koeffizienten derjenigen
Potenzreihen, welche eine schlichte Abbildung des
Einheitskreises vennitteln', Sitzungsber. Preuss. Akad. Wiss.
Phys-Math. Kl. (1916),940-955.
[2] GOLUZIN, G.M.: Geometric theory of functions of a complex
variable, Trans!. Math. Monographs, 26, Amer. Math. Soc.,
1969 (translated from the Russian).
[3] MILIN, I.M.: Univalent functions and orthonormal systems,
Trans!. Math. Monographs, 49, Amer. Math. Soc., 1977
(translated from the Russian).
[4] HAYMAN, W.K.: 'The asymptotic behaviour of p-valent func-
tions', Proc. London Math. Soc. (3) 5 (1955),257-284.
[5A] OZAWA, M.: 'On the Bieberbach conjecture for the sixth coeffi-
cient', Kodai Math. Sem. Rep. 21 (1969), 97-128.
[5B] OZAWA, M.: 'An elementary proof of the Bieberbach conjec-
ture for the sixth coefficient', Kodai Math. Sem. Rep. 21 (1969),
129-132.
[6] PEDERSON, R.N. and SCHIFFER, M.M.: 'A proof of the Bieber-
bach conjecture for the fifth coefficient', Arch. Rat. Mech. and
Anal. 45 (1972), 161-193.
[7] FITZGERALD, C.H.: 'Quadratic inequalities and coefficient esti-
mates for the fifth coefficient', Arch. Rat. Mech. and Anal. 46
(1972), 356-368.
[8] SHIROKOV, N.A.: 'On a regularity theorem of Hayman', J.
Soviet Math. 2, no. 6 (1974), 693-710. (Zap. Nauchn. Sem.
Leningr. Otd. Mat. Inst. Steklov. 24 (1972), 182-200)
[9] BAZILEVICH. I.E.: Mathematics in the USSR during 40 years:
1917-1957, Vo!' I, Moscow, 1959, pp. 444-472.
[10] HOROWITZ, D.: 'A refinement eStimate for univalent func-
tions', Proc. A mer. Math. Soc. 54 (1976), 176-178.
E. G. Goluzina
Actually Bieberbach proved
1 C2 1.;;;2 and then asked in a footnote if perhaps generally
1 c n I.;;;n, ct. [1].
In the laborious progress on the Bieberbach conjecture
from 1950 until 1975, M. Schiffer's variational methods
played a major role. For a detailed history, see Duren's
book [A3].
The case n=4 was proved by P.R. Garabedian and M.M.
Schiffer in 1955.
In 1984 the Bieberbach conjecture was established in
complete generality by the French-born U.S. mathematician
Louis de Branges [A 1], [A2]. In fact, he proved a remark-
able hypothesis of the Soviet mathematicians N.A. Lebedev
and I.M. Milin which was known to be even stronger than
the Bieberbach conjecture, cf. [3], [A3]. The Lebedev-Milin
hypothesis asserted that for n = 2, 3,
QnCO) = ~1 {k 1 Yk 12 -1-}cn
k=1 k
- k) .;;; O. C)
Here the numbers Yk are the logarithmic coefficients of fCz) ,
385
BIEBERBACH CONJECTURE
defined by
Every function fez) in S is the starting point fez, 0) of a
Loewner parametric family {f(z, t)}, t~O, of univalent func-
tions. The latter map the disc onto a continuously-increasing
family of domains with limit C. Normalized so that
fez, t) / e t is in S, the functions of a Loewner family satisfy
Loewner's partial differential equation
af af
at = Z~p(z, t),
where Rep(z, tO [A4] , [A3]. De Branges introduced
functionals Qn(t) corresponding to the logarithmic coeffi-
cients Yk(t) for fez, t) / e t , replacing the weight factors n - k
in ("') by undetermined weights 0nk(t) with initial value n - k.
By ingenious use of Loewner's differential equation he could
show that the weights 0nk(t) may be chosen such that
Q~(t)~O for all t~O and all n. The proof also used a sophis-
ticated positivity result for hypergeometric functions due to
R. Askey and G. Gasper. Since Qn(t)~O as t~oo it followed
that Qn(O)~O. The latter inequality established the
Lebedev - Milin hypothesis ("') and thereby also
Bieberbach's conjecture I en I ~n. The case of equality was
dealt with both by de Branges [A2] and by C.H. FitzGerald
and C. Pommerenke [A5].
References
[A 1] BRANGES, L. DE: 'A proof of the Bieberbach conjecture' .
Preprint E-5-84. Steklov Math. Inst. Leningrad (1984), 1-21.
[A2] BRANGES, L. DE: 'A proof of the Bieberbach conjecture', Acta
Math. 154 (1985),137-152.
[A3] DUREN, P.L.: Univalent functions, Springer, 1983.
[A4] POMMERENKE, c.: Univalent functions, Vandenhoeck and
Ruprecht, 1975.
[A5] FITZGERALD, C.H. and POMMERENKE, c.: 'The de Branges
theorem on univalent functions', Trans. Amer. Math. Soc.
290 (1985), 683-690.
AMS 1980 Subject Classification: 30C50
BIEBERBACH - ElLENBERG FUNCfIONS in the
disc I z I < 1 - The class R of functions j (z), regular in
the disc I z I <1, which have an expansion of the form
j(Z)=CIZ + ... +cnz n+ ...
and which satisfy the condition
*'
j(Zdj(Z2) 1, 1 Zl 1 <1, 1 Z2
This class of functions is a natural extension of the
class B of functions j (z), regular in the disc I z I < 1,
with an expansion (1) and such that I j (z) I < 1 for
I z I < 1. The class of univalept functions (cf. Univalent
function) in R is denoted by R. The functions in R were
named _after L. Bieberbach [1], who showed that for
j(z)ER the inequality
1 CI 1 ~ 1
is valid, while equality holds only for the function
j(z)=eiOz, where () is real, and after S. Eilenberg [2],
who proved that the inequality (2) is valid for the
whole class R. It was shown by W. Rogosinski [3] that
every function in R is subordinate (cf. Subordination
386
principle) to some function in R. Inequality (2) yields
the following sharp inequality for j(z)ER:
1 j' (z) 1 ~ Ill.:=- ~(i ~ I, 1 z 1 < 1. (3)
The following bound on the modulus of a function in
R has been obtained: If j(z)ER, then
1 fez) 1 ~ (1_:2)1/2' 1 z 1 =r, O<r<I, (4)
and (4) becomes an equality only for the functions
+j(zeiO;r), where () is real and
. _ (l-r 2 )1/2 z
f(z,r)- l' .
+lrz
The method of the extremal metric (cf. Extremal metric,
method of the) provided the solution of the problem of
tpe maximum and rnipmum of I j (z) I in the class
R(c) of functions in R with a fixed value I Cl L=c,
O<c.:;;;l, in the expansion (1): For j(z)ER(c),
O<c < 1, the following sharp inequalities are valid:
ImH(ir;r, c) ~ 1 f(re i9 ) 1 ~ ImF(ir;r, c). (5)
Here the functions w=H(z;r, c) and w=F(z;r, c)
map the disc I z I < 1 onto domains which are sym-
metric with respect to the imaginary axis of the w-
plane, and the boundaries of which belong to the union
of the closures of certain trajectories or orthogonal tra-
jectories of a quadratic differential in the w-plane with a
certain symmetry in the distribution of the zeros and
Eoles [4], [5]. Certain optimal results for functions in
R(c) were obtained by the simultaneous use of the
method of the extremal metric and the symmetrization
method [4].
_ Many results obtained for the functions in the classes
Rand R are consequences of corresponding results for
systems of functions mapping the disc I z I < 1 onto
disjoint domains [6]. The analogue of R for a finitely-
connected domain G without isolated boundary points
) and not containing the point z = 00, is the class
(1)
Ra(G),aEG, of functions j(z) regular in G and satisfy-
ing the conditions j(a)=O, j(ZI)j(Z2)=I=l, where
1 <1. z I, Z 2 are arbitrary points in G. The class Ra( G)
extends the class Ba( G) of functions j (z), regular in G
and such that j(a)=O, I j(z) I <1 in G. The following
sharp estimate is an extension of the result of
Bieberbach - Eilenberg and of inequality (3) to func-
tions of class Ra(G): If j(Z)ERa(G), then
I/(z) 1 ~ 1 1- j2(Z) 1 F'(z, z), ZEG.
where F(z, b), bEG, is that function in Bb(G) for
(2) which F'(b, b)=max I feb) I in this class.
References
[IJ BIEBERBACH, L.: 'Ueber einige Extremalprobleme im Gebiete
der konformen Abbildung', Math. Ann. 77 (1916), 153-172.
[2J ElLENBERG, S.: 'Sur quelques proprietes topologiques de la sur-
face de sphere', Fund. Math. 25 (1935),267-272.
[3J ROGOSINSKI, W.: 'On a theorem of Bieberbach - Eilenberg', 1.

London Math. Soc. (1) 14 (1939), 4-11.
[4] JENKINS, J.A.: Univalent functions and conformal mappings,
Springer, 1958.
[5] JENKINS, J.A.: 'On Bieberbach - Eilenberg functions III',
Trans. Amer. Math. Soc. 119, no. 2 (1965), 195-215.
[6] LEBEDEV, N.A.: The area principle in the theory of univalent
functions, Moscow, 1975 (in Russian).
G. V. Kuz'mina
Editorial comments.
References
[A 1] DUREN, P.L.: Univalent functions, Springer, 1983.
AMS 1980 Subject Classification: 30C70
BIEBERBACH POLYNOMIALS - Extremal polyno-
mials which approximate a function that conformally
maps a given simply-connected domain onto a disc.
These polynomials were first studied by L. Bieberbach
[1] in the context of the problem of the approximate
computation of a conformal mapping.
Let G be a simply-connected domain in the finite
part of the plane bounded by a curve r, and let the
function w = cp(z) map this domain conformally and
univalently onto the disc I w I <ro under the condi-
tions cp(zo)=O and Ij>'(zo) = 1, where Zo is an arbitrary
fixed point of G and ro depends on zoo The polynomial
'lTn(z) which minimizes the integral
I(Fn) = f f G 1 F~(z) 12 dx dy
in the class of all polynomials Fn(z) of degree n subject
to the conditions Fn(zo)=O and F~(zo)= I is called the
Bieberbach polynomial. In the class of all functions
which are analytic in the domain G and which satisfy
the same conditions, this integral is minimized by the
mapping function w =cp(z). If the contour r is a Jordan
curve, the sequence {'ITn(z)} converges uniformly to the
function cp(z) inside G. In the closed domain G there
need not be convergence [2]. If the contour r satisfies
certain additional smoothness conditions, the sequence
{'lTn(z)} converges uniformly in the closed domain, and
the rate of convergence depends on the degree of
smoothness of r.
References
[I] BIEBERBACH, L.: 'Zur Theorie und Praxis der konfonnen
Abbildung', Rend Circ. Mat. Palermo 38 (1914),98-112.
[2] KEWYSH, M.V.: 'Sur l'approximation en moyenne quadratique
des fonctions analytiques', Mat. Sb. 5 (47), no. 2 (1939), 391-
401.
[3] MERGELYAN, S.N.: Some questions of the constructive theory of
functions, Moscow, 1951 (in Russian).
[4J SUETIN, P.K.: 'Polynomials orthogonal over a region and
Bieberbach polynomials', Proc. Steklov Inst. Math. 100 (1974).
(Trudy Mat. Inst. Steklov. 100 (1971
P.K. Suetin
Editorial comments. A good additional reference is [A 1].
References
[A1] GAlER, D.: Vorlesungen uber Approximation im Komplexen,
Birkhauser, 1980.
AMS 1980 Subject Classification: 30C10, 30C70
BIFURCA nON
BIFACfORIAL MAPPING - A mapping f of a topo-
logical space X into a topological space Y, in which for
any covering of the inverse image [-I (y) of any point
Y E f (X) by sets open in X it is possible to select a fin-
ite number of sets so that y is located inside the image
of their union. It is particularly important that the pro-
duct of any collection of bifactorial mappings is a
bifactorial mapping. Bifactorial mappings constitute an
extensive class of factorial mappings (cf. Factorial
mapping), but nevertheless preserve the fine topological
properties of spaces. Thus, continuous bifactorial s-
mappings preserve a pointwise-countable base, and a
factorial s-mapping of a space with a pointwise-
countable base onto a space of pointwise-countable
type is bifactorial.
v. V. Filippov
AMS 1980 Subject Classification: 54C10
BIFUNCfOR - A mapping T: 2f X SB~C, defined on
the Cartesian product of two categories 2f and SE with
values in C, which assigns to each pair of objects A E 2f,
B ESE some object C EC and to each pair of mor-
phisms
a:A ~ A', fJ:B ~ B'
the morphism
T(a, fJ): T(A', B) -> T(A, B'). (I)
The following conditions
T(lA, 1B ) = 1T(A. B),
T(a' oa, fJ' of3) = T(a, fJ')oT(a', fJ), (2)
must also be met. In such a case one says that the func-
tor T is contravariant with respect to the first argument
and covariant with respect to the second.
V. E. Govorov
Editorial comments. What is described above is a
bifunctor contravariant in its first argument and covariant in
its second. A bifunctor covariant in both arguments, the
more fundamental notion ([A 1D, has (1) and (2) replaced by
T(a, f3): T(A, B) ~ T(A, B), (1 ')
T(a'a, fJfJ) = T(a, fJ")oT(a, fJ) (2')
Similarly one can define bifunctors contravariant in both
arguments and covariant in the first and contravariant in the
second argument.
References
[A 1] MITCHELL, B.: Theory of categories, Acad. Press, 1965,
AMS 1980 Subject Classification: 18A22
BIFURCATION - A term in certain branches of
mathematics, applied to situations in which some object
:I) = :I)(A) depends on a parameter A (which is not
necessarily scalar) and which is such that in any neigh-
bourhood of a certain value Ao of that parameter (a
bifurcation value or a bifurcation point) the qualitative
properties under consideration of the object :I)(A) are
387
BIFURCATION
not the same for all X. The corresponding rigorous
definitions vary in different cases, but mainly follow (in
a more or less modified form) the two variants below:
A) The qualitative properties of the object ~ being
studied are the existence of other objects 0, associated
with it in some manner. The bifurcation is character-
ized by the fact that when X is varied, the objects 0
appear or disappear (in particular, they may coincide
or else one object may generate several ones). See below
(Paragraph 1.
B) The first step is to determine under what cir-
cumstances two objects ~(X) are considered to be
equivalent. (This definition must be such that all the
qualitative properties of interest be identical for
equivalent objects.) A change in the qualitative proper-
ties of ~(X) in a neighbourhood of a bifurcation point
means, by definition, that X-values with non-equivalent
~(X) are found near this point. See below (Paragraph
2.
I) In the theory of operators the initial object ~(X) is
a non-linear operator fl>(x, X) in a real Banach space,
depending on a real parameter X, defined in a neigh-
bourhood of the point x =0 and such that fl>(0, X)-O.
For each fixed X one associates new objects 0 to this
fl>; they are the solutions x of the non-linear operator
equation fl>(x, X)=x. A bifurcation point is a point at
which a new, non-trivial, solution of this equation is
generated. It is in fact a point Ao such that for any t:>0
there exists a X, I X- Ao I <t:, for which the equation
x = fl>(x, X) has a solution x(X) satisfying the conditions
0< " x(X) "<t:. If fl>(x, X) AAx, where A is a linear
completely-continuous operator, the concept of a bifur-
cation point coincides with the concept of a characteris-
tic value of A.
If fl>(x, X) is a non-linear completely-continuous
operator which is continuously Frechet differentiable
and is such that fl>x(O, X) M, then only the charac-
teristic values of A can be bifurcation poir) s of fl>. It
was found by a topological method [1], [2] that each
odd-fold (in particular, simple) characteristic value of A
is a bifurcation point of fl>. An analogous sufficient
condition for the case of even-fold characteristic values
is formulated using the concept of the rotation of a vec-
tor field.
If x =0 is a non-isolated solution of the equation
x =fl>(x, ,\), then Ao is a bifurcation point of fl>. It was
shown by a variational method [1], [2] that if fl>(x) is a
non-linear completely-continuous operator in a Hilbert
space which is the gradient of a weakly continuous
functional, while A = fl>' (0) is a completely-continuous
self-adjoint operator, then all characteristic values of A
are bifurcation points of fl>. The concept of a bifurca-
tion point is also modified in the case of large solutions
X('\)-HX) for X~Ao. The importance of these concepts
388
and results consists in the fact that, subject to relatively
weak limitations, the branching of the solution x =0
can be established; in particular, it is possible to prove
that the solution of the non-linear problem is not
unique. The methods of the theory of branching of
solutions of non-linear equations often yield more accu-
rate information [5].
2) The theory of smooth dynamical systems studies
one-parameter (and sometimes also two-parameter [6])
families of flows (and cascades; only the former will be
considered here), and the conditions under which the
bifurcation is 'typical', i.e. preserves its character under
a small change of the family in question [9]; both vari-
ants A) and B) above are useful. In the second variant
two flows are considered equivalent if there exists a
homeomorphism of the phase space which converts the
trajectories of the one into trajectories of the other
while preserving the direction of motion. There exists a
completely satisfactory theory of bifurcations of one-
parameter families of flows with a two-dimensional
phase space [7], [9], and a local variant referring to a
neighbourhood of an equilibrium position or of a
periodic solution in the n-dimensional case [6].
In variant A) the objects 0 under study, which are
associated to the given dynamical system, are the posi-
tions of an equilibrium or periodic solution, and occa-
sionally certain invariant manifolds (mainly tori) or
hyperbolic invariant sets. The 'growth' of these objects
- taking place 'locally' near an equilibrium position,
or near a periodic solution, or 'semi-locally' in a neigh-
bourhood of a 'closed contour' formed by several tra-
jectories which tend to an equilibrium position or to a
periodic solution - is considered. A case of bifurca-
tion which is connected with a similar contour in a cer-
tain sense, but which takes place (with a change of the
parameter ,\) prior to its arrival, is also possible [8]. The
growth of periodic solutions is often conveniently con-
sidered by rewriting the differential equation and the
periodicity condition in the form of an integral equa-
tion to which suitable methods are applied [5].
3) Various bifurcations of different objects (both the
initial objects and those associated to them) are
encountered in the theory of singularities of mappings.
As a result, this term (or rather a term derived from it)
is used in several different ways [10], [6], [11], but it is
more common to assign independent names to the
corresponding concepts. They include, for instance, ver-
sal families (or deformations) [6], [11], [12], which
describe, in a certain sense, all possible bifurcations
which may take place under a small deformation of the
object under consideration. These include, in particu-
lar, seven elementary catastrophes [12], which are
represented by 'typical' k-parametric (k ";;4) families of
functions which include a function with a degenerate

critical point and which are defined in a neighbourhood
of that point; they accordingly describe the correspond-
ing bifurcation. In non-Soviet literature on the theory
of singularities, the term 'catastrophes' is often
employed for 'bifurcations'.
References
[1) KRASNOSEL'sKIi, M.A.: Topological metods in the theory ofnon-
linear integral equations, Pergamon, 1964 (translated from the
Russian).
[2) KREiN, S. (ED.): Functional analysis, Handbook of Math.
Libraries, Moscow, 1964 (in Russian).
[3) KRAsNOSEL'sKIi, M.A: Positive solutions of operator equations,
Wolters-Noordhoff, 1964 (translated from the Russian).
[4) VAiNBERG, M.M.: Variational methods for the study of nonlinear
operators, Holden-Day, 1964 (translated from the Russian).
[5) VAiNBERG, M.M. and TRENOGIN, V.A.: Theory of branching of
solutions of non-linear equations, Noordhoff, 1974 (translated
from the Russian).
[6) ARNOL'D, V.I.: 'Lectures on bifurcations and versal families',
Russian Math. Surveys 27, no. 5 (1972), 54-123. (Uspekhi Mat.
Nauk 27, no. 5 (1972), 119-lS4)
[7) ANDRONOV, A.A., LEONTOVICH, E.A., GoRDON, 1.1. and
MAIER, A.G.: Theory of bifurcations of dynamic systems on a
plane, Israel Progr. Sci. Transl., 1971 (translated from the Rus-
sian).
[SA) GAVRlLOV, N.K. and SHIL'NIKOV, L.P.: 'On three-dimensional
systems close to systems with a structurally unstable homoc-
linic curve 1', Math. USSR-Sb. 17, no. 4 (1972), 467-4S5.
(Mat. Sb. 88 (130), no. 4 (1972),475-492)
[SB) GAVRlLOV, N.K. and SHIL'NIKOV, L.P.: 'On three-dimensional
systems close to systems with a structurally unstable homoc-
linic curve II', Math. USSR-Sb. 19, no. 1 (1973), 139-156.
(Mat. Sb. 90 (132), no. 1 (1973), 139-156)
[9) PEIXOTO, M.M.: 'On bifurcations of dynamical systems', in
Proc. Internat. Congress Mathematicians Vancouver, 1974, Vol.
2, Vancouver, 1974, pp. 315-319.
[10) THOM, R.: 'The bifurcation subset of a space of maps', in
Manifolds, Amsterdam 1970, Springer, 1971, pp. 202-20S.
[II) ARNOL'D, V.I.: 'Critical points of smooth functions and their
normal form', Russian Math. Surveys 30, no. 5 (1975), 1-75.
(Uspekhi Mat. Nauk 30, no. 5 (1975), 3-65)
[I 2) BROCKER, P. and LANDER, L.: Differentiable germs and catas-
trophes, Cambridge Univ. Press, 1975.
D. V. Anosov
V.A. Trenogin
Editorial comments. A standard reference for bifurcation
theory is [A 1]. The presence of symmetry in a bifurcation
problem can often be a powerful tool [A2]. For more pre-
cise details in the case of bifurcations of solutions to equa-
tions cf. Branching of solutions.
References
[A1) CHOW, S.-N. and HALE, J.K.: Methods of bifurcation theory,
Springer, 1982.
[A2) SALLINGER, D.H.: Branching in the presence of symmetry,
SIAM,1983.
AMS 1980 Subject Classification: 35832, 58E07, 58-
XX
BIHARMONIC FUNCTION A function
u (x) = u (x I, . . . ,xn ) of real variables, defined in a
domain D of the Euclidean space Rn, n ;"2, with con-
tinuous partial derivatives up to the fourth order
inclusive, that satisfies in D the equation
BIHARMONIC FUNCTION
il 2 u _ il(ilu) = 0,
where d is the Laplace operator. This equation is
known as the biharmonic equation. The class of bihar-
monic functions includes the class of harmonic func-
tions and is a subclass of the class of poly-harmonic
functions (cf. Harmonic function; Poly-harmonic func-
tion). Each biharmonic function is an analytic function
of the coordinates Xi'
From the point of view of practical applications the
biharmonic functions in two variables u(xJ, X2) are the
most important. Such biharmonic functions can be
represented using harmonic functions UI, U2 or VI, V2,
in the form
or
u (x J, X2) = (r2 - ra)v 1(x J, X2)+ V2(X 1, X2),
where r2 =xT + x~, while r6 is a constant. The main
boundary value problem for biharmonic functions is
the following: To find a biharmonic function in the
domain D that is continuous together with its first-
order derivatives in the closed domain D= Due and
satisfies on the boundary C the conditions
ule = !1(S), ~~ Ie = h(s), (*)
where au / an is the derivative with respect to the nor-
mal on C, and 11 (s), h (s) are given continuous func-
tions of the arc length s on the contour C. The above
representations of biharmonic functions give the solu-
tion of problem (*) in an explicit form for the case of
the disc D, starting from the Poisson integral for har-
monic functions [1].
Biharmonic functions in two variables may also be
represented as follows:
u(XJ, X2) = Re{zcj>(z)+x(z)} =
= ~ {zcj>(z)+zcj>(z)+x(z)+x(z)}, Z=Xl -iX2,
using two analytic functions cp(z), X(z) of the complex
variable z = x I + ix 2. This representation makes it pos-
sible to reduce the boundary value problem (*) for an
arbitrary domain D to a system of boundary value
problems for analytic functions, for which a method of
solving was developed in detail by G.V. Kolosov and
N.!. Muskhelishvili. This method was developed while
solving various planar problems of elasticity theory (cf.
Elasticity theory, planar problem of), in which the main
biharmonic function is the stress function or the Airy
function [2], [3] (d. Airy functions).
References
[I] TlCHONOFF, AN. [AN. TIKHoNov) and SAMARSKll, A.A.: Dif-
ferentialgleichungen der mathematischen Physik, Deutsch. Ver-
lag Wissenschaft., 1959 (translated from the Russian).
(2) MUSKHELISHVILI, N.I.: Some basic problems of the mathemati-
389
BIHARMONIC FUNCTION
cal theory of elasticity, Noordhoff, 1975 (translated from the
Russian).
[3] LAVRENT'EV, M.A. and SHABAT, B.V.: Methoden der komplexen
Funktionentheorie, Deutsch. Verlag Wissenschaft., 1967
(translated from the Russian).
E.D. Solomentsev
Editorial comments. An axiomatic treatment of bihar-
monic functions, similar to that of harmonic functions, is
given in [A1], [A2].
References
[A 1] SMYRNELIS, E.P.: 'Axiomatique des fonctions biharmoniques',
Ann. Insf. Fourier (Grenoble) 25, no. 1 (1975), 35-97.
[A2] SMYRNELIS, E.P.: 'Axiomatique des fonctions biharmoniques',
Ann. Inst. Fourier (Grenoble) 26, no. 3 (1976), 1-47.
AMS 1980 Subject Classification: 31 A30, 31830,
31C05
BmOLOMORPHIC MAPPING, holomorphic isomor-
phism, holomorphism, pseudo-conformal mapping - A
generalization of the concept of a univalent confonnal
mapping to the case of several complex variables. A
holomorphic mapping of a domain D ccn onto a
domain D' ccn is said to be a biholomorphic mapping if
it is one-to-one. A biholomorphic mapping is non-
degenerate in D; its inverse mapping is also a biholo-
morphic mapping.
A domain of holomorphy is mapped into a domain of
holomorphy under a biholomorphic mapping; holo-
morphic, pluriharmonic and plurisubharmonic func-
tions are also invariant under a biholomorphic map-
ping. If n > 1, biholomorphic mappings are not confor-
mal (except for a number of linear mappings) and the
Riemann theorem is invalid for biholomorphic map-
pings (e.g. a ball and a polydisc in C2 cannot be biho-
lomorphically mapped onto each other). A biholo-
morphic mapping of a domain D onto itself is said to
be a (holomorphic) automorphism; if n > 1, there exist
simply-connected domains without automorphisms
other than the identity mapping.
References
[I] SHABAT, B.V.: Introduction to complex analysis, 1-2, Moscow,
1976 (in Russian).
E.D. Solomentsev
Editorial comments. Concerning boundary behaviour of
biholomorphic mappings the following results have been
obtained. C. Fetterman's theorem: A biholomorphic mapping
between strongly pseudo-convex domains with Coo -smooth
boundary extends Coo -smoothly to a diffeomorphism
between the closures of the domains, see [A3]. The same
result holds if the domains are only pseudo-convex and one
of them satisfies condition R for the Bergman projection,
see [A2]. For strongly pseudo-convex domains with C k _
boundary, k>2, C k - 1 -, extendability was obtained (>0 if
k=2, 3, ... , =0 otherwise) by L. Lempert and by S.
Pincuk. For (weakly) pseudo-convex domains with real-
analytic boundary one has even holomorphic extension to a
neighbourhood of the closure, see [A 1]. Similar results were
obtained for proper holomorphic mappings.
390
A biholomorphic mapping 'is proper (i.e. the pre-image
of a compact set is compact), since ,-1
is continuous.
Riemann's theorem does not hold in the following sense:
There is no proper holomorphic mapping from the polydisc
in en onto the ball in em for any n, m> 1, cf. [A4]. Thus,
function theory in en, n;;;.1, is strongly related to the
domain of definition of the functions. For function theory in
the (unit) ball of en see [A5]; for function theory in
polydiscs see [A6]. For entire holomorphic mappings and
their value distribution see [A7].
References
[A1] BAOUENDI, M.S., JACOBOWITZ, H. and TREVES, F.: 'On the
analyticity of CR mappings', Ann. ot Math. 122 (1985), 365-
400.
a
[A2] BELL, ST.: 'Biholomorphic mappings and the problem', Ann.
otMath. 114 (1981),103-113.
[A3] FEFFERMAN, c.: 'The Bergman kernel and biholomorphic
mappings of pseudoconvex domains', Inv. Math. 26 (1974),
1-65.
[A4] KRANTZ, S.G.: Function theory ot several complex variables,
Wiley, 1982, Chapt. 10.
[A5] RUDIN, W.: Function theory in the unit ball otcn, Springer,
1980.
[A6] RUDIN, W.: Function theory in polydiscs, Benjamin, 1969.
[A7] GRIFFITHS, PH.A.: Entire holomorphic mappings in one and
several variables, Princeton Univ. Press, 1976.
AMS 1980 Subject Classification: 32HXX
BIJECTION, bijective mapping, of a set A into a set
B - A mapping f: A --"?B under which different elements
in A have different images in f (A) = B. In other
words, f is a one-to-one mapping of A onto B, i.e. a
mapping that is both an injection and a surjection. A
bijection establishes a one-to-one correspondence
between the elements of the sets A and B. A bijection
of a set A onto itself is also called a permutation of A.
o.A. Ivanova
AMS 1980 Subject Classification: 04A05
BIUNEAR DIFFERENTIAL - An analytic differen-
tial on a Riemann surface, depending on two points P
and Q, and having the form
fez, ndz dt
where z and ~ are local uniformizing parameters III a
neighbourhood of P and Q respectively, and f (z, n is
an analytic function of z and r Bilinear differentials
are used to express many functionals on a finite
Riemann surface.
References
[I] SCHIFFER,M. and SPENCER, D.C.: Functionais offinite
Riemann surfaces. Princeton Univ. Press. 1954.
E.D. Solomentsev
AMS 1980 Subject Classification: 30F30
BIUNEAR FORM on a product of modules V X W -
A bilinear mapping f: V X W --"?A, where V is a left uni-
tary A-module, W is a right unitary A-module, and A is

a ring with a unit element, which is also regarded as an
(A, A)-bimodule. If V= W, one says that fis a bilinear
form on the module V, and also that V has a metric
structure given by f Definitions involving bilinear
mappings make sense also for bilinear forms. Thus, one
speaks of the matrix of a bilinear form with respect to
chosen bases in V and W, of the orthogonality of ele-
ments and submodules with respect to bilinear forms,
of orthogonal direct sums, of non-degeneracy, etc. For
instance, if A is a field and V = W is a finite-
dimensional vector space over A wi th basis e], ... , en,
then for the vectors
v =V I eI + ... + l'n en
and
w=wlel + ... +wne"
the value of the form will be
n
1(1', w) = ~ a,lv,It'I'
/)~ I ring). See also Quadratic form for the structure of bil-
where ai,. = f (e;. e). The polynomial 2: I,)n . _ ai, V; w) in
th e vanables v], ... , vn , W] . . . 'W n IS sometImes
-].
identified with f and is called a bilinear form on V. If
the ring A is commutative, a bilinear form is a special
case of a sesquilinear fonn (with the identity automor-
phism).
Let A be a commutative ring. A bilinear form on an
A -module V is said to be symmetric (or anti-symmetric
or skew-s)immctric) if for all VI, V 2 E V one has
f(V],V2)=j(V2,vd (or f(V],V2)=-f(V2,vd), and is
said to be alternating if f(v, v)=O. An alternating bil-
inear form is anti-symmetric; the converse is true only
if for any aEA it follows from 2a=0 that a=O. If V
has a finite basis. symmetric (or anti-symmetric or
alternating) forms on V and only such forms have a
symmetric (anti-symmetric, alternating) matrix in this
basis. The orthogonality relation with respect to a sym-
metric or anti-symmetric form on V is symmetric.
A bilinear form f on V is said to be isometric with a
bilinear form g on W if there exists an isomorphism of
A-modules : V ~ W such that
g(Cv), Cw =I (v, w)
for all v E V. This isomorphism is called an isometry of
the form and, if V = Wand f = g, a metric automorphism
of the module V (or an automorphism of the form f)
The metric automorphisms of a module form a group
(the group of automorphisms of the form f); examples of
such groups are the orthogonal group or the symplectic
group.
Let A be a skew-field and let f be a bilinear form on
V X W; let the spaces V / W 1- and W / V ~ be finite-
dimensional over A; one then has
dim V / W ~ = dim W / V ~ ,
and this number is called the rank of
BILINEAR FORM
dimensional and f is non-degenerate, then
dim V = dim W
and for each basis v], ... ,Vn in V there exists a basis
w], ... ,Wn in W which is dual with respect to f; it is
defined by the conditionf(v;, w)=8 ij , where 8ij are the
Kronecker symbols. If, in addition, V= W, then the
submodules V ~ and W ~ are said to be the right and
the left kernel of f, respectively; for symmetric and
anti-symmetric forms the right and left kernels are
identical and are simply referred to as the kernel.
Let f be a symmetric or an anti-symmetric bilinear
form on V. An element v E V for which f(v, v)=O is
said to be an isotropic element; a submodule MeV is
said to be isotropic if M n
M ~ =1= {O}, and totally isotro-
pic if Me M ~. Totaliy isotropic sub modules play an
important role in the study of the structure of bilinear
forms (d. Witt decomposition; Witt theorem; Witt
inear forms.
.
Let A be commutative, let HomA (V, W) be the A-
module of all A-linear mappings from V into W, and
let L 2 (V, W, A) be the A-module of all bilinear forms
on VX W. For every bilinear form f on VX Wand for
each Vo E V, the formula
't"o(w) = 1(1'0, W), WE W,
defines an A-linear form on W. Correspondingly, for
Wo E W the formula
rr,wo(v) = 1(1', wo). VEV,
defines an A-linear form on V. The mapping
l( vo~lf,vo is an element of
HomAV, HomA(W, A,
The mapping rf in
HomAW, HomA(V ,A
is defined in a similar way. The mappings f~/f and
f~rf define isomorphisms between the A-modules
L2(V, W, A) -> Hom4(V, HomA(W, A
and
L2(V, W, A) -> HomAW, Hom4(V, A.
A bilinear form f is said to be left-non-singular (respec-
tively, right-non-singular) if If (respectively, rf) is an iso-
morphism; if f is both left- and right-non-singular, it is
said to be non-singular; otherwise it is said to be singu-
lar, A non-degenerate bilinear form may be singular.
For free modules V and W of the same finite dimen-
sion a bilinear form f on V X W is non-singular if and
only if the determinant of the matrix of f with respect
to any bases in V and W is an invertible element of the
ring A. The following isomorphisms
I
f If V is finite- Hom A (V. V) --> L 2 (V, W, A)
391
BILINEAR FORM

and r
A =B and A is commutative, then the set L 2(V, W, H)
HOmA(W, W) ~ L 2 (V, W, A), of all bilinear mappings VX W~H is an A-module
with respect to the pointwise defined operations of
given by a non-singular bilinear form f, are defined by
addition and multiplic~tion with elements in A, while
the formulas
the correspondence f ~f establishes a canonical isomor-
1(q,Xv, w) = f(q,(v), w)
phism between the A-module L 2(V, W, H) and the A-
and
r(~Xv, w) = f(v, ~w.
module L (V w, H) of all linear mappings from
VWinto H.
The endomorphisms epE HomA (V, V) and Let V and W be free modules with bases Vi, i EI, and
~EHomA(W, W) are said to be cOnjugate with respect to Wi' j EJ, respectively. A bilinear mapping f is fully
the form f if ~ = (r -1 0/)(ep).
determined by specifying f(Vi, Wi) for all i EI, j EJ,
References since for any finite subsets I' CI, f CJ, the following
[1] BOURBAKI, N.: Elements of mathematics. Algebra: Algebraic
structures. Linear algebra, 1, Addison-Wesley, 1974,
formula is valid:
f( ~.a;v;, ~.wA] ~ aJ(v;, wj)bj.
Chapt.l;2 (translated from the French).
= (*)
[2] LANG, S.: Algebra, Addison-Wesley, 1974. ieI jeJ ier,
[3] ARTIN, E.: Geometric algebra, Interscience, 1957. jeJ

v.L. Popov
AMS 1980 Subject Classification: 15A63
BILINEAR FUNCfIONAL - A bilinear mapping from
a module over a commutative ring K (a K-module) into
K itself (considered as a K-module).
M.l. Voitsekhovskii
AMS 1980 Subject Classification: 16A64
BILINEAR INTEGRAL FORM - The double integral
b b
J(q" ~)= j jK(x, s)</>(x~s)dx ds,
a a
where K(x, s) is a given (usually complex-valued)
square-integrable function of real variables, and cJ>(x),
1[;(x) are arbitrary (also complex-valued) square-
integrable functions, while 1[;(s) is the complex conju-
gate function of 1[;(s). If 1[;(s) = cJ>(s), J (ep, ep) is said to be
a quadratic integral form.
B. V. Khvedelidze
AMS 1980 Subject Classification: 46CXX
BILINEAR MAPPING, bilinear function - A mapping
ffrom the product VX W of a left unitary A-module V
and of right unitary B-module W into an (A, B)-
bimodule H, satisfying the conditions
f(v +v', w) = f(v, w)+ f(v', w);
f(v, w+w') = f(v, w)+f(v, w');
f(av, w) = af(v, w);
f(v, wb) = f(v, w)b;
where v, v' E V, W, W' E W, a EA, b EB are arbitrarily
Conversely, after the elements hij EH, i EI, j EJ, have
been chosen arbitrarily, formula (*), where
f(Vi, Wi)=hij' defines a bilinear mapping from VX W
into H. If I and J are finite, the matrix Ilf(Vi' wi)11 is
said to be the matrix of f with respect to the given
bases.
Let a bilinear mapping f: VX W ~H be given. Two
elements V E V, W E W are said to be orthogonal with
respect to f if f (u, w) = O. Two subsets X C V and
YeW are said to be orthogonal with respect to f if any
x EX is orthogonal to any y E Y. If X is a submodule in
V, then
= {wEW:f(x,w)=O for allXEX},
X-L
which is a submodule of W, is called the orthogonal
submodule or the orthogonal complement to X. The
orthogonal eomplement Y ~ of the submodule Y in W
is defined in a similar way. The mapping f is said to be
right-degenerate (left-degenerate) if V ~=i={O}
(W~=i={O}). The sub modules V~ and W~ are called,
respectively, the left and right kernels of the bilinear
mapping f If V~ = {O} and W~ = {O}, then f is said
to be non-degenerate; otherwise it is said to be degen-
erate. The mapping f is said to be a zero mapping if
V ~ = W and W ~ = V.
Let v., i EI, be a set of left A-modules, let Wi, i EI,
be a set of right B-modules, let fi be a bilinear mapping
from Vi X W; into H, let V be the direct sum of the A-
modules v.,
and let W be the direct sum of the B-
modules Wi' The mapping f: VX W ~H, defined by
the rule
f(~v;, ~w;] = (v;, Wi),
lEI lEI lEI

chosen elements, and A and B are rings with a unit ele-
ment. The tensor product V W over Z has the
natural structure of an (A, B)-bimodule. Let
q,: VX W~VWbe a canonical mapping; any bilinear
mapping f will !hen induce a homomorphism of
(A, B)-bimodules f: V W ~H for which f= foep. If
is a bilinear mapping and is said to be the direct sum of
the mappings Ji. This is an orthogonal sum, i.e. the sub-
module Vi is orthogonal to the sub module Jfj with
respect to fif i=i=j.
The bilinear mapping f is non-degenerate if and only
if fi is non-degenerate for all i EI. Moreover, if f is

392

non-degenerate then one has
Vjl. = ~ nJ, Wjl. = ~~.
Joj:-j Joj:-j
If A = B = H, a bilinear mapping is called a bilinear
form.
References
[I) BOURBAKI, N.: Elements of mathematics. Algebra: Algebraic
structures. Linear algebra, I, Addison-Wesley, 1974,
Chapt.!;2 (translated from the French).
[2) LANG, S.: Algebra, Addison-Wesley, 1974.
v.L. Popov
AMS 1980 Subject Classification: 16A64, 15A69,
13CXX
BIMATRIX GAME - A finite non-cooperative game
between two players. A bimatrix game is specified by
two matrices A = II aU II and B = I bi) I of the same
dimension m X n, which are, respectively, the payoff
matrices (or gain matrices) of the players I and II. The
strategy of player I is the selection of a row, that of
player II the selection of a column. If player I chooses i
(l';:;;i';:;;m), while player II chooses j (l';:;;j.;:;;n), their
respective payoffs (gains) will be ail and bi); if
aij + bi) = 0 for all i, j, the bimatrix game becomes a
matrix game. The theory of bimatrix games is the sim-
plest branch of the general theory of non-cooperative
games, but even bimatrix games are not always solvable
according to Nash or strongly solvable. Various algo-
rithms are available to find equilibrium solutions in
bimatrix games: a method of describing the subma-
trices of A, B yielding all extreme points of the set of
equilibrium solutions [1], [2]; and methods which
reduce the problem of finding the equilibrium solutions
of a birnatrix game to a problem of quadratic program-
ming [3], [4], [5].
References
[I] VOROB'EV, N.N.: 'Equilibrium points in bimatrix games',
Theory Probab. Appl3, no. 3 (1958), 297-309. (Teor. Veroyat
nost. i Primenen. 3, no. 3 (1958),318-331)
[2] KUHN, H.W.: 'An algorithm for equilibrium points in bimatrix
games', Proc. Nat. A cad. Sci. USA 47 (1961), 1657-1662.
[3] MIllS, H.: 'Equilibrium points in infinite games', J. Soc. Ind.
Appl. Math. 8, no. 2 (1960), 397-402.
[4] MANGASARIAN, O.L.: 'Equilibrium points of bimatrix games',
J. Soc. Ind Appl. Math. 12, no. 4 (1964), 778-780.
[5] LEMKE, C.E. and HOWSON, JR., J.T.: 'Equilibrium points of
bimatrix games', J. Soc. Ind Appl. Math. 12, no. 2 (1964),
413-423.
E.B. Yanovskaya
AMS 1980 Subject Classification: 90005
BIMODAL DISTRIBUTION, doubly-peaked distribu-
tion - A probability distribution which is characterized
by the fact that the probability curve has two local
maxima, corresponding to two values of the modes (cf.
Mode). Such a distribution is often the result of 'mix-
ing' two normal distributions (cf. Normal distribution).
BINARY COMPUTING SYSTEM
See also Multimodal distribution; Unimodal distribution.
A. V. Prokhorov
AMS 1980 Subject Classification: 62E99
BIMODULE, double module - An Abelian group B
that is a left module over a ring R and a right module
over a ring S, and is such that (rb)s = r (bs) for all r E R,
b EB, s ES. One says that this is the situation RBs , or
that B is an (R, S)-birnodule. The birnodule B may be
regarded as a left R S' -module, where S' is the ring
which is dually isomorphic (anti-isomorphic) to S,
while denotes the tensor product over the ring of
integers, and (r s)b = rbs. For every left R -module M
one has the situation RME , where E is the ring of
endomorphisms of M. Any ring A can be given the
natural structure of an (A, A )-birnodule.
L.A. Skornyakov
Editorial comments. A bimodule morphism is a mapping
from a bimodule RBS into a bimodule ReS that is left R-linear
and right 5-linear. The category of (R, S)-bimodules with
bimodule morphisms is a Grothendieck category.
The centre of an (R, R)-bimodule (also called an R-
bimodule) RBR is defined to be the set
ZR(B) = {xEB: rb=brfor all rER}. Clearly ZR(B) is a two-
sided ZR(R)-module.
AMS 1980 Subject Classification: 16A64
BIMORPIDSM, bijective morphism in a category -
One of the categoric-theoretical generalizations of the
concept of a bijective mapping between sets. A mor-
phism u in a category C is said to be a birnorphism if it
is both a monomorphism and an epimorphism in C. A
product of birnorphisms is a birnorphism, i.e. the
birnorphisms form a subcategory containing all isomor-
phisms. In the category of sets and the category of
groups every birnorphism is an isomorphism. However,
the categories of rings, topological spaces, or Abelian
groups without torsion contain birnorphisms that are
not isomorphisms.
1. V. Dolgachev
M. S. Tsalenko
AMS 1980 Subject Classification: 18A20
BINARY COMPUTING SYSTEM - A positional sys-
tem for the representation of numbers in base 2 (cf.
Numbers, representations of). In such a system the
arithmetic operations can be very simply executed. For
instance, the usual addition and multiplication with
single-valued numbers are reduced in the system to:
I + 1 = 10; 11 = 1.
V. 1. N echaev
Editorial comments. The term binary arithmetic is also
used.
AMS 1980 Subject Classification: 68M05, 10AXX
393
BINARY FORM
BINARY FORM - A form in two variables, viz. a
homogeneous polynomial
where the coefficients at. k =0, ... ,n, belong to a
given commutative ring with a unit element. Such a ring
may be the ring Z of integers, the ring of integers of
some algebraic number field, the field R of real
numbers or the field C of complex numbers. The
number n is called the degree of the form. If n = 2, f is
called a binary quadratic fonn.
The theory of forms includes algebraic. (theory of
invariants), arithmetic (representation of numbers by
forms) and geometric (theory of arithmetical minima of
forms) approaches. The purpose of the algebraic theory
of binary forms (in R or q is to construct a complete
system of invariants of such forms under linear
transformations of variables with coefficients of the
same field (cf. Invariants, theory of; see also [2], Chapt.
5). The arithmetic theory of binary forms studies
Diophantine equations of the form
j(x,y) = b,
where aQ, ... ,an, b EZ, their solvability and their solu-
tions in the ring Z. The most important result is Thue's
theorem and its generalizations and sharpenings (d.
Thue-Siegel-Roth theorem). See [5], Chapts. 9 - 17,
and the MordeD conjecture on the solvability of such
equations in the field Q and the possible number of
solutions. The theory of arithmetical minima of binary
forms is part of the geometry of numbers. The arithmet-
ical minimum of a form f is defined as the quantity
m(j) = (x,y)eiW\(O,O) I j(x, y) I,
It has been proved for the case n = 3 that
{ II DD /2311/4
/49 11/ if D >0,
4
m(j)";; if D<O,
where D is the discriminant of f, which, in the present
case, IS
18aoa1a2a3 +aTa~ -4ooa~ -403a1 -27a6aj.
These estimates cannot be improved.
References
[1] BOREVICH, Z.I, and SHAFAREVICH, I.R.: Number theory, Acad.
Press, 1966 (translated from the Russian). German transla-
tion: Birkhauser, 1966.
[2] GUREVICH, G.B.: FoundatioTL'J of algebraic invariant theory,
Moscow-Leningrad, 1948 (in Russian).
[3] LANDAU, E. and WALFISZ, A.: Diophantische Gleichungen mit
endlich vielen LOsungen, Deutsch. Verlag Wissenschaft., 1959.
[4] LEKKERKERKER, e.G.: Geometry of numbers, Wolters-
Noordhoff, 1969.
[5] MORDELL, L.J.: Diophantine equatioTL'J, Academic Press, 1969.
A. V Malyshev
AMS 1980 Subject Classification: 10CXX
394
BINARY LIE ALGEBRA, BL-algebra - A linear alge-
bra A over a field F any two elements of which gen-
erate a Lie sub algebra. The class of all binary Lie alge-
bras over a given field F generates a variety which, if
the characteristic of F is different from 2, is given by
the system of identities
x 2 = J(xy, x,y) = 0, (*)
where
J(x,y, z) = (xy)z +(yz)x + (zx)y.
If the characteristic of F is 2 and its cardinal number is
not less than 4, the class of binary Lie algebras can be
defined not only by the system of identities (*), but
also by the identity
J([(~)y]x, x,y) = o.
The tangent algebra of an analytic local alternative loop
is a binary Lie algebra and vice versa.
References
[I] MAL'TSEV, A.I.: 'Analytic loops', Mat. Sb. 36 (78), no. 3
(1955), 569-575 (in Russian).
[2] GAiNOV, A.T.: 'Binary Lie algebras of characteristic two', Alge-
bra and Logic 8, no. 5 (1969), 287-297. (Algebra j Logika 8,
no. 5 (1969), 505-522) T ~
A . . Gamov
AMS 1980 Subject Classification: 17B99
BINARY P-ADIC GROUP - An infinite group G of
square matrices of order 2
where a, b, c, d are integral p-adic numbers (cf. p-adic
number) satisfying the following conditions:
ad-bc = 1, c _ 0 (modp), d _ 1 (modp).
The quotient groups of such groups of the form G / N,
where N is the n-th member of the lower central series
of G or the n-th term of the derived series (the series of
higher commutators of G), are examples of finite p-
groups having certain extremal properties.
A.I. Kostrikin
Editorial comments.
References
[A1] HUPPERT, B.: Endliche Gruppen, 1, Springer, 1979.
AMS 1980 Subject Classification: 20C99, 20E99
BINARY QUADRATIC FORM - A quadratic form in
two variables, i.e. having the form
j = j(x,y) = ax 2+bxy+ry2. (*)
If a, b and c are integers, the binary quadratic form is
said to be integral. The expression d = ac - b 2 /4 is
called the discriminant or determinant of the binary qua-
dratic form. The expression b 2 -4ac is also sometimes
referred to as the discriminant. The arithmetic theory of
binary quadratic forms originated with P. Fermat, who

proved that any prime number of the form 4k + 1 can
be represented as the sum of two squares of integers.
The theory of quadratic forms was completed by J.L.
Lagrange and by C.F. Gauss. This theory is a special
case of the theory of quadratic forms in n variables; the
arithmetic theory of binary quadratic forms is
equivalent to the theory of ideals in quadratic fields,
and is one of the origins of algebraic number theory
(cf. Quadratic form; Quadratic field).
The number of genera of binary quadratic forms
with discriminant d equals 2S -J, where s is the number
of different prime divisors of d, except for d .1
(mod4), d-O (mod8), when s is increased by one; if
- d is a square, the number of different binary qua-
dratic forms is doubled. The number r(d, m) of essen-
tially different primitive representations o~ a number. m
by a complete system of binary quadratic fo~s WIth
discriminant d is equal to the number of solutlOns of
the equation X2 -d (modm).
As in the general case, there exists an algorithm which
reduces the problem of solving a given second-degree
Diophantine equation in two unknowns (in particu~ar,
an equation f(x,y)=m) to the problem on the anth-
metic equivalence of two binary quadratic forms.
All integral automorphisms of a primitive form f
with a=FO can be represented in the form
Ilt-~:/2 t+~:/211'
where t 2 +du 2 = 1, and 2t and u are integers (cf. PeU
equation). Therefore, the problem on the equivalence of
two forms is solved by the reduction theory of binary
quadratic forms. The reduction theory of positive-
definite binary quadratic forms is a special case of the
reduction theory of positive-definite quadratic forms
according to H. Minkowski. The reduction theory of
integral indefinite binary quadratic forms can be
reduced to the reduction theory of quadratic irrationali-
ties ([2], [3]).
An important role in the theory of numbers is played
by the arithmetic function h (d) - the number of
classes of primitive integral binary quadratic forms with
determinant d. It is known that h (d)< + 00. Some idea
of the rate of increase of the function h (d) can be
obtained from Siegel's theorem: Let d>O, then for any
>0 there exist constants c. c:
and >0 which satisfy
the condition
c;d 1 / 2 -, < h(d) < c,d 1 /2+,
(a similar formula is also valid for d<O).
Let I::. be an integer, 1::._1 or 0 (mod 4), sUEPose that
if s211::. then s=l or s=2, and let F=Q(v'I::.) be the
quadratic field which is obtained by adjoining Vi to
the field of rational numbers. A correspondence has
BINARY RELATION
been established between the integral ideals [aJ, a2l of
the field F and the integral quadratic forms
N(alx+a7)')
f(x,y) = N[a], a2l
with - A/ 4 as determinant. This results in a one-to-
one correspondence (up to a transition to conjugate
classes of ideals) between the ideal classes of the field F
and the classes of binary quadratic forms. In this
correspondence, multiplication of ideal classes defines a
composition of classes of binary quadratic forms.
As in the case of forms in n variables, the theory of
binary quadratic forms can be generalized to include
forms (*) with coefficients a, b and c in a given alge-
braic number field.
There are various variants in the definitions of an
integral form, the discriminant of the form, equivalence
of forms, and the classes and the genera of forms. The
definition of integral forms given above is due to L.
Kronecker. Gauss [1 1 stipulated that b be even. In
determining the equivalence (and the class of forms),
only the substitutions with discriminant + 1 may be
considered; in other cases discriminants + 1 are con-
sidered. The definition of a genus given in [6] is wider
than that given by Gauss.
References
[1] GAUSS, C.F.: Disquisitiones arithmeticae, Yale Univ. Press,
1966 (translated from the Latin).
[2] VENKOV, B.A.: Elementary number theory, Wolters-Noordhoff,
1970 (translated from the Russian).
[3] JONES, B.W.: The arithmetic theory of quadratic forms, Math.
Assoc. Amer., 1950.
[4] GEL'FOND, A.O. and LINNIK, Yu.V.: Elementary methods in the
analytic theory of numbers, M.I.T., 1966 (translated from the
Russian).
[5] LANDAU, E.: Vorlesungen iiber Zahlentheorie, 3, Hirzel, 1927.
[6] BOREVICH, Z.1. and SHAFAREVICH, 1.R.: Number theory, Acad.
Press, 1966 (translated from the Russian). German transla-
tion: Birkhauser, 1966.
[7] O'MEARA, O.T.: Introduction to quadratic forms, Springer, 1973.
A. V. Malyshev
AMS 1980 Subject Classification: 10CXX
BINARY RELATION - A two-place predicate on a
given set. The term is sometimes used to denote a sub-
set of the set A XA of ordered pairs (a, b) of elements
of a given set A. A binary relation is a special case of a
relation. Let RCA XA. If (a, b)ER, then one says that
the element a is in binary relation R to the element b.
An alternative notation for (a, b )ER is aRb.
The empty subset 0 in A XA and the set A XA itself
are called, respectively, the nil relation and the universal
relation in the set A. The diagonal of the set A XA, i.e.
the set 1::.= {(a, a): a EA}, is the equality relation or the
identity binary relation in A.
Let R, R 1, R 2 be binary relations in a set A. In
addition to the set-theoretic operations of union
395
BINARY RELATION

Rl UR2' intersection Rl nR2, and complementation (2)

R' = (A XA) \ R, one has the inversion
R- 1 = {(a, b): (b, a)ER}, Binomial coefficients, as well as the arithmetical trian-
gle, were known concepts to the mathematicians of
as well as the operation of multiplication: antiquity, in more or less developed forms. B. Pascal
R 1oR 2 = {(a, b): (3CEA)(aR1c and cR 2 b)}. (1665) conducted a detailed st\ldy of binomial coeffi-
The binary relation R - 1 is said to be the inverse of R. cients. The binomial coefficients are also connected by
Multiplication of binary relations is associative, but as many useful relationships other than (2), for example:
a rule not commutative.
A binary relation R in A is said to be I) reflexive if
ACR; 2) transitive if RoR CR; 3) symmetric if [~J k~JkJ [~=kJ, n";;;m";;;N-n;
R- l CR; and 4) anti-symmetric if R nR~1 CA. If a
binary relation has some of the properties 1), 2), 3) or f(-I)kk m
k=O
[ZJ = 0, m=O, ... ,N-I;
4), the inverse relation R -1 has these properties as
(3)
well. The binary relation RCA XA is said to be func-
tional if R -loR CA.
k~Jkr = [2nnJ;

The most important types of binary relations are

equivalences, order relations (total and partial), and
k~JZJk(k-I)'" (k-r+l) =

functional relations (d. Equivalence; Order relation; = N(N-I) (N-r+lpN-r;

Functional relation).
D.M. Smirnov k~O(-li [ZJk(k -I)' .. (k-r+ I) = O.
AMS 1980 Subject Classification: 04A05
In particular, (3) yields
BINARY UNIT in information theory - A synonym
for bit.
AMS 1980 Subject Classification: 94A17, 68MOS (4)

BINOMIAL - A sum or a difference of two algebraic

expressions, said to be the terms of the binomial, e.g.
a +b, 5x _y2 / (1 + x 2 ), etc. For powers of a binomial, The use of the Stirling formula yields approximate
i.e. expressions of the type (x +y t, see Newton bino- expressions for binomial coefficients. Thus, if N is
mial. much larger than n:

AMS 1980 Subject Classification: 12-XX Nn

[~J ~-"
n.
BINOMIAL COEFFICIENTS - The coefficients at the In the case of a complex number a, binomial coeffi-
powers of z in the decomposition of the Newton bino- cients are generalized according to the formula
t.
mial (1 + z The binomial coefficients are denoted by

[~J
[aJn = a(a-l) (a-n+l)
n." n>,
o [aJ
0
= I.
or by CRr, and are given by In this generalization, some of the relations (2) - (4) are
en - N!
preserved, but usually in a modified form. Thus,
[NJ
n -
-
N - n!(N-n)!
(I)

N(N-I)'" (N-n+l)
, ' O";;;n";;;N. [~J+ [n~d = [~!lJ;
n.
The first-mentioned notation is due to L. Euler; the ~ [k J = 2u, Re a> - I;
k=O
notation CRr appeared in the nineteenth century, and is
probably connected with the interpretation of the bino- k~(_I)k [k J = 0, Rea>O.
mial coefficients C'J., as the number of distinguishable
non-ordered combinations (d. Combination) from N For tables of binomial coefficients see [2], [3].
different objects with n objects in each combination. References
Binomial coefficients are most conveniently written as [I) KORN, G.A. and KORN, T.M.: Mathematical handbook for

numbers in the arithmetical triangle, or Pascal triangle, scientists and engineers, McGraw-Hill, 1968.
[2) BOL'sHEv, L.N. and SMIRNOV, N.V.: Tables of mathematical
the construction of which is based on the following statistics, Libr. of mathematical tables, 46, Nauka, Moscow,
property of binomial coefficients: 1968 (in Russian). Processed by L.S. Bark and E.S. Kedova.

396
 BINOMIAL SERIES

[3) Table of binomial coefficients, Cambridge, 1954. Moivre- Laplace theorem):

E.D. Solomentsev
F(y) = ell [y-n p +O.5]+R (y )
AMS 1980 Subject Classification: 05A 10 Vnp(l-p) n ,p,
where
BINOMIAL DISTRIBUTION, Bernoulli distribution - Rn(y,p) = O(n-l/2)
The probability distribution of a random variable X
uniformly for all real y. There also exist other, higher
which assumes integral values x =0, ... ,n with the
order, normal approximations of the binomial distribu-
probabilities
tion.
P{X=x} = bAn,p) = [~)px(l_p)n-x, If the number of independent trials n is large, while
the probability p is small, the individual probabilities
where G) is the binomial coefficient, and p is a parame-
bAn, p) can be approximately expressed in terms of the
ter of the binomial distribution, called the probability
Poisson distribution:
of a positive outcome, which can take values in the
interval O~ ~ 1. The binomial distribution is one of
the fundamental probability distributions connected
bAn,p) = [~)px(l-pr-x ~ (;r e- np .
with a sequence of independent trials. Let Y\, Y 2 , .. , If n-HXJ and O<C''';;y~C (where c and C are con-
be a sequence of independent random variables, each stants), the asymptotic formula
one of which may assume only one of the values I and
F(y)= ~AX,e-A+O(n-2),
o with respective probabilities p and 1 - P (i.e. all Yj are x=ox.
binomially distributed with n = I). The values of Y j where A= (2n - [y])p / (2 - p), is uniformly valid with
may be treated as the results of independent trials, with respect to all p in the interval O<p < 1.
Y j = 1 if the result of the i-th trial is 'positive' and The multinomial distribution is the multi-dimensional
Y j =0 if it is 'negative'. If the total number of indepen- generalization of the binomial distribution.
dent trials n is fixed, such a scheme is known as Ber-
noulli trials, and the total number of positive results, References
[I) GNEDENKO, B.V.: The theory ofprobability, Chelsea, reprint,
x = Y + . . . + Yn ,
I n;;;.l, 1962 (translated from the Russian).
[2) FELLER, W.: An introduction to probability theory and its appli-
is then binomially distributed with parameter p. cations, 1-2, Wiley, 1957-1971.
The mathematical expectation Ezx (the generating [3] PROHOROV, Yu.V. [Yu.V. PROKHOROV) and ROZANov, YU.A.:
Probability theory, Springer, 1969 (translated from the Rus-
function of the binomial distribution) for any value of z
sian).
is the polynomial fpz + (1-p)r, the representation of [4) PROKHORov, Yu.V.: 'Asymptotic behaviour of the binomial
which by Newton's binomial series has the form distribution', Selected Translations in Math. Stat. and Probab. 1
(1961), Amer. Math. Soc. (Uspekhi Mat. Nauk 8, no. 3 (1953),
b o+b1z + ... +bnz n. 135-142)
[5) BOL'SHEV, L.N. and SMIRNOV, N.V.: Tables of mathematical
(Hence the very name 'binomial distribution'.) The
statistics, Libr. of mathematical tables, 46, Nauka, Moscow,
moments (cf. Moment) of a binomial distribution are 1965 (in Russian). Processed by L.S. Bark and E.S. Kedova.
given by the formulas L.N. Bol'shev
AMS 1980 Subject Classification: 62EXX
= np,EX
OX = E(X -np)2 = np(l-p), BINOMIAL SERIES - A power series of the form
E(X -np)3 = np(l-pXI-2p).
n~o [~)zn = 1+ [l)z+ [i)Z2+ ... ,
The binomial distribution function is defined, for any
real y, O<y <n, by the formula where n is an integer and 0: is an arbitrary fixed
number (in general, a complex number), z = x + iy is a
F(y) = P{X";;;u} = x~J~)px(l_p)n-x, complex variable, and the
where [y 1is the integer part of y, and

F(y) = 1
I
jtlY](I-t)n-lY]-1 dt,
B([y]+I,n-[y]) p
are the binomial coefficients. For an integer 0: = m ;;;;;. 0
the binomial series reduces to a finite sum of m + 1
B(a, b) is Euler's beta-function, and the integral on the terms
right-hand side is known as the incomplete betajunction. ( l+z)m = l+mz+ m(m -1)z2+ ... +zm
As n-HX), the binomial distribution function is 2! '
expressed in terms of the standard normal distribution which is known as the Newton binomial. For other
function ell by the asymptotic formula (the de values of 0: the binomial series converges absolutely for

397
BINOMIAL SERIES
I z I < I and diverges for I z I > 1. At points of the unit
circle I z I = I the binomial series behaves as follows:
I) if Rea>O, it converges absolutely at all points; 2) if
Reao;;;;; -I, it diverges at all points; and 3) if
-I <Reao;;;;;O, the binomial series diverges at the point
z = - I and converges conditionally at all other points.
At all points of convergence, the binomial series
represents the principal value of the function (I + z t
which is equal to one at z = O. The binomial series is a
special case of a hypergeometric series.
If z =x and a are real numbers, and a is not a non-
negative integer, the binomial series behaves as follows:
I) if a >0, it converges absolutely on - I";;;x 0;;;;; 1; 2) if
a 0;;;;; -I, it converges absolutely in -I <x< 1 and
diverges at all other values of x; and 3) if -I <a 0;;;;;0,
the binomial series converges absolutely in - I <x < I,
converges conditionally at x = I, and diverges for
x = - 1; for I x I > I the binomial series always
diverges.
Binomial series were probably first mentioned by I.
Newton in 1664 - 1665. An exhaustive study of bino-
mial series was conducted by N.H. Abel [I], and was
the starting point of the theory of complex power
series.
References
[I] ABEL,N.H.: 'Untersuchungen tiber die Reihe
I+m /2x+m(m-I)/2'lx 2 + ... ',J. Reine Angew. Math. 1
(1826), 311-339.
[2] KNopp, K.: Theorie und Anwendung der unendlichen Reihen,
Springer, 1964. English translation: Blackie, 1951.
[3] MARKUSHEVICH, A.I.: Theory offunctions of a complex variable,
I, Chelsea, reprint, 1977 (translated from the Russian).
E.D. Solomentsev
AMS 1980 Subject Classification: 40A05, 30810
BINORMAL - The straight line passing through a
point M 0 of a curve L perpendicular to the osculating
plane to L at Mo. If r=r(t) is a parametrization of L,
then the vector equation of the binormal at M 0
corresponding to the value to of the parameter t has the
form
E. V. Shikin
Editorial comments. This definition holds for space
curves for which r"Uo) does not depend linearly on r'Uo),
i.e. the curvature should not vanish.
For curves in a higher-dimensional Euclidean space, the
binormal is generated by the second normal vector in the
Frenet frame (cf. Frenet trihedron), which is perpendicular
to the plane spanned by r'Uo) and r"Uo) and depends
linearly on r'Uo), r"Uo), r"'Uo) (cf. [A1]).
References
[A1] SPIVAK, M.: Differentia/geometry, 2, Publish or Perish, 1970.
AMS 1980 Subject Classification: 53A04
BIORTIIOGONAL SYSTEM - A pair of sets {at}
398
and {~,}, t E T, of elements of a (topological) vector
space X and the dual (topological) space X*, respec-
tively, which satisfies the conditions
~(as) = <~(' as> = 0
if t=l=s, and =1=0 if t = s (here, <.,. > is the canonical
bilinear form coupling X and X*). For instance, a
biorthogonal system consists of a Schauder basis and
the set formed by the expansion coefficients of x in it.
In a Hilbert space H with scalar product <','> and
basis {at} the set {bs } satisfying the condition
<a(,bs>=lisl>
where ~st = I if t =s and ~st =0 if t=l=s, is also a basis; it
is said to be the basis dual to {at} and, since H = H* ,
the sets {at} and {bt } form a biorthogonal system. In
particular, a basis in H is said to be orthonormal if its
dual to itself.
However, there also exist biorthogonal systems which
do not even form a weak basis; an example is the set of
functions eikX, k EZ, x ER, in the space of continuous
periodic functions with the norm II f II = sup If (x) I
M.I. Voftsekhovskif
AMS 1980 Subject Classification: 46A20
BIPLANAR SPACE - A real (2n + I)-dimensional
projective space, with two non-intersecting n-
dimensional subspaces which are real (biplanar spaces
of hyperbolic type) or complex conjugate (biplanar
spaces of elliptic type), the fundamental group of which
consists of the projective transformations that map each
of these subspaces into itself. These two n-dimensional
subspaces are called the absolute planes. The linear
congruence of real straight lines intersecting both abso-
lute planes is said to be an absolute congruence. This
congruence serves as a real model of an n-dimensional
projective space over the algebra of double or complex
numbers. If n = I, a biplanar space is said to be a bi-
axial space. The study of properties of geometrical fig-
ures in biplanar spaces which are preserved under the
operation of the fundamental group is the subject of
hiplanar geometry. The bi-axial geometry in which one
studies the theory of curves, surfaces and complexes of
straight lines, has been investigated in detail.
A.P. Shirokov
AMS 1980 Subject Classification: 53A35, 51 N25
BIPOLAR COORDINATES - The numbers l' and a
which are connected with the Cartesian orthogonal
coordinates x and y by the formulas
asinh'T a sin a
x = cosh1'-cosa' Y = cosh'T-cosa'
where O';;;;a<7T, -00<1'<00. The coordinate lines are
two families of circles (1'=const) with poles A and B
 BlRATIONAL GEOMETRY

and the circles orthogonal with these (a=const). BIQUADRATIC EQUATION - An equation of the
(Joe=: const type
a.x 4 +bx 2 +c = 0,
where a, h, c are given complex numbers and a*O. The
substitution y =x 2 converts a biquadratic equation into
a quadratic equation.
AMS 1980 Subject Classification: OOA99, 12E12

BIRATIONAL GEOMETRY - The branch of alge-

braic geometry in which the main problem is the clas-
The Lame coefficients are: sification of algebraic varieties up to birational
equivalence (cf. Birational mapping). Over a fixed field
of constants k every class of birationally equivalent
LT = La = (COshT-cosai'
varieties defines a finitely generated field over k, which
The Laplace operator is: is isomorphic to the field of rational functions on any
variety in this class. Conversely, to each such field
!Jot = -(COshT-cosai
1 32 r 32 r J.
[ ~+~
a2 3~ 3~ corresponds a class of birationally equivalent varieties
- the model(s) of this field. Thus, the birational clas-
Bipolar coordinates in space (bispherical coordinates)
sification of algebraic varieties is equivalent to the clas-
are the numbers a, T and </>, which are connected with
sification (up to a k-isomorphism) of the finitely gen-
the orthogonal Cartesian coordinates x, y and z by the
erated fields regular over k.
formulas:
x = a sin a cos </> The most general birational invariant is the dimension
cosh T - cos a ' of an algebraic variety. For one-dimensional algebraic
a sin a sin</> varieties - irreducible algebraic curve~ - each class of
Y = COshT-cosa'
birational equivalence contains a non-singular model -
a SinhT
Z = -----'---
cosh cos a'
T-
a smooth projective curve, which is unique up to k-
isomorphisms. Accordingly, the birational classifica-
where - 00 <a< 00, O:S;;;T<'lT, O:S;;;</><2'lT. The coordi- tion of algebraic curves is reduced to the classification,
nate surfaces are spheres (a=const), the surfaces up to k-isomorphisms, of smooth projective curves,
obtained by rotation of arcs of circles (T=const), and which leads to the moduli problem. For dimensions ;;;;.2
half-planes passing through the Oz-axis. The system of the problem becomes much more complicated. The very
bipolar coordinates in space is formed by rotating the existence of a smooth model constitutes a problem of
system of bipolar coordinates on the plane Oxy around resolution of singularities of algebraic varieties, which at
the Oz-axis. the time of writing (1986) has been positively solved
The Lame coefficients are: only for surfaces and for varieties of arbitrary dimen-
sion over a field of characteristic zero. In such a case, if
La = LT = (COshT-cosa)2' such models exist at all, there are infinitely many of
L -
a 2 sin2 a
--=--==:........:....---:- them in the class of birationally equivalent varieties.
</> - (COShT-cosoi The minimal models (cf. Minimal model) take a special
The Laplace operator is: place among such models. Their birational classification

!Joj = (cosbr-cosha)3
a 2 sino
[..!.3T [ sin a
cosh T- cos a 3T
l1.J +
is often identical with classification up to a k-
isomorphism, as in the case of curves. However, this is

n]
not true in the general case, even for (both rational and
+..!.3a [ sino
cosh T- cos a 3a
l1J + 1
sin a( cosh T- cos a) 3cfl .
ruled) surfaces.
The principal results in the classification of algebraic
surfaces were obtained by geometers of the Italian
References
[I] MADELUNG, E.: Die mathematischen Hilfsmittel des Physikers,
school [1]. Only isolated results are so far (1986) avail-
Springer, 1957. able [3], [7], [8] for varieties of dimensions of ;;;;'3.
D.D. Sokolov The principal discrete birational invariants of smooth
Editorial comments. complete algebraic varieties over a field k of charac-
References teristic zero include the arithmetic genus, the geometric
[A 1] VEBLEN, O. and WHITEHEAD, 1.H.c.: The foundations of dif- genus, the multiple genus, the dimension of the space
ferential geometry, Cambridge Univ. Press, 1932.
of regular differential forms, the Severi torsion, the fun-
AMS 1980 Subject Classification: 51 N 10 damental group and the Brauer group.

399
BlRATIONAL GEOMETRY
One of the most important problems in birational
geometry is the problem of rationality of algebraic
varieties, i.e. the problem of the description of rational
varieties - varieties which are birationally equivalent
to a projective space, cf. Rational variety.
If the field of constants is algebraically non-closed,
the problems in birational geometry are closely
connected with the arithmetic of algebraic varieties (cf.
Algebraic varieties, arithmetic of). The important prob-
lem in such a case is that of describing the birational
k-forms on a given variety V over a field k, in particu-
lar if, say, V=PZ is a projective space over k [2]. An
important part of this problem is the description of the
group of birational transformations of the variety V.
References
[1) Algebraic surfaces, Moscow, 1965 (in Russian).
(2) MANIN, YU.N.: Cubic forms. Algebra, geometry, arithmetic,
North-Holland, 1974 (translated from the Russian).
(3) ROTH, L.: Algebraic threefolds, Springer, 1955.
(4) SHAFAREVICH, LR.: Basic algebraiC geometry, Springer, 1977
(translated from the Russian).
(5) BALDASSARRI, M.: Algebraic varieties, Springer, 1956.
(6) DoLGACHEV, LV. and ISKovsKIKH, V.A.: 'The geometry of
algebraic varieties', J. Soviet Math. 5, no. 6 (1976), 803-864.
(lto~ Nauk. i Tekhn. Algebra Topol. Geom. 12 (1974),
77-170)
(7) J. Soviet Math. 13, no. 6 (1980). (Itogi Nauk. i Tekhn. Sovrem.
Prob!. Mat. 12 (1979
(8) IflAKE, S.: AlgebraiC geometry, Springer, 1982.
(9) HARTSHORNE, R.: Algebraic geometry, Springer, 1977.
I. V. Dolgachev
V.A. Iskovskikh
Editorial comments. A field extension K / k is regular if
k is algebraically closed in K and K is linearly disjoint from
the algebraic closure k of k over k. If k is algebraically
closed, then every extension of k is regular and there is an
arrow-reversing equivalence of categories between the
category of varieties and dominant rational mappings over k
and the category of finitely generated field extensions, cf.
[8], Sect. 104. For dimension 3 many new results have
recently become available in the (birational) classification of
algebraic varieties, cf. [A 1].
References
[A 1] CONTE, A. (ED.): Algebraic threefolds, Leet. Notes in Math,
Vol. 947, Springer, 1982.
AMS 1980 Subject Classification: 14E05
BIRATIONAL MAPPING, birational isomorphism - A
rational mapping between algebraic varieties inducing
an isomorphism of their fields of rational functions. In
a more general setting, a rational mapping of schemes
f : X ~ Y is said to be a birational mapping if it satisfies
one of the following equivalent conditions: 1) there
exist dense open sets U C X and V C Y such that f is
defined on U and realizes an isomorphism of sub-
schemes flu: U~V; 2) if {Xd;d, (Yj}jd are the sets
of generic points of the irreducible components of the
schemes X and Y respectively, f induces a bijective
400
correspondence between the sets a: I ~J and an iso-
morphism of local rings (J) X, Xi ~ (J) Y'Y"i) for each i EI.
If the schemes X and Y are irreducible and reduced,
the local rings of their generic points become identical
with the fields of rational functions on X and Y,
respectively. In such a case the birational mapping
f : X ~ Y induces, in accordance with condition 2), an
isomorphism of the fields of rational functions:
R(Y}~:~R(X).
Two schemes X and Y are said to be birationally
equivalent or birationally isomorphic if a birational map-
ping f : X ~ Y exists. A birational morphism is a special
case of a birational mapping.
The simplest birational mapping is a monoidal
transformation with a non-singular centre. For smooth
complete varieties of dimension .,;;;2 any birational
mapping may be represented as the composite of such
transformations and their inverses. At the time of writ-
ing (1986) this question remains open in the general
case.
References
[1) SHAFAREVICH, I.R.: Basic algebraiC geometry, Springer, 1977
(translated from the Russian).
(2) HARTSHORNE, R.: Algebraic geometry, Springer, 1977.
I. V. Dolgachev
V.A. Iskovskikh
AMS 1980 Subject Classification: 14E05
BIRATIONAL MORPIDSM - A morphism of
schemes that is a birational mapping. The most impor-
tant examples of birational morphisms include: a nor-
malization, a blowing up and a monoidal transforma-
tion. Any proper birational transformation between
regular two-dimensional schemes can be decomposed
into monoidal transformations with non-singular cen-
tres [2] (cf. Monoidal transformation). This is not true
in dimensions higher than two.
References
[1) GROTHENDIECK, A. and DIEUDONNE, J.: 'Elements de
geometrie algebrique', Publ. Math. IHES 8 (1960).
(2) SHAFAREVICH, I.R.: Lectures on minimal models and birational
transformations of two-dimensional schemes, Tata Inst. Fundam.
Res., 1966.
(3) SHAFAREVICH, I.R.: Basic algebraic geometry, Springer, 1977
(translated from the Russian).
(4) HARTSHORNE, R.: Algebraic geometry, Springer, 1977.
I. V. Dolgachev
V.A. Iskovskikh
AMS 1980 Subject Classification: 14E05
BIRATIONAL TRANSFORMATION - A birational
mapping of an algebraic variety (or scheme) into itself.
Also called sometimes a birational automorphism. The
group of all birational transformations of an algebraic
variety is canonically isomorphic to the group of auto-

morphisms of its field of rational functions over the
field of constants. Examples of birational transforma-
tions include the Cremona transformations (d.
Cremona transformation), in particular the standard
quadratic transformation of a projective plane, given by
the formula
(xo, x" X2) ~ (X]X2, XOX2, XOx]),
where (x 0, X I , X 2) are homogeneous coordinates in the
projective plane.
References
(1) ISKOVSKIKH, V.A.: 'Birational automorphisms of three-
dimensional algebraic varieties', J. Soviet Math. 13, no. 6
(1960),815-868. (ltogi Nauk. i Tekhn. Sovrem. Probl. Mat. 12
(1979), 159-239)
l. V Dolgachev
VA. Iskovskikh
AMS 1980 Subject Classification: 14E05
BIRKHOFF ERGODIC THEOREM - One of the
most important theorems in ergodic theory. For an
endomorphism T of a space X with a a-finite measure /L
Birkhoffs ergodic theorem states that for any function
jELI(X, /L) the limit
1 n -] - associative algebra R by replacing the operation of
lim - "2.f(T k x) = f(x)
n----'+oo n k =0
(the time average or the average along a trajectory)
exists _almost everywhere (for almost all x EX). More-
over,jELI(X, /L), and if p.{X)<oo, then
Ifd/L = [ldp.
For a measurable flow {T/} in the space X with a a-
finite measure /L Birkhoffs ergodic theorem states that
for any functionjELI(X, /L) the limit
1T _
lim -Tjf(T,x)dt = f(x)
0
T---H()
exists almost everywhere, with the same properties of f
Birkhoffs theorem was stated and proved by G.D.
Birkhoff [1]. It was then modified and generalized in
various ways (there are theorems which contain, in
addition to Birkhoffs theorem, also a number of state-
ments of a somewhat different kind which are known
in probability theory as ergodic theorems (d. Ergodic
theorem); there also exist ergodic theorems for more
general semi-groups of transformations [2]). Birkhoffs
ergodic theorem and its generalizations are known as
individual ergodic theorems, since they deal with the
existence of averages along almost each individual tra-
jectory, as distinct from statistical ergodic theorems -
the von Neumann ergodic theorem and its generaliza-
tions. (In non-Soviet literature the term 'pointwise
ergodic theorem' is often used to stress the fact that the
averages are almost-everywhere convergent.)
BIRKHOFF - WITT THEOREM
References
[I) BIRKHOFF, G.D.: 'Proof of the ergodic theorem', Proc. Nat.
Acad. Sci. USA 17 (1931),656-660.
(2) KATOK, A.B., SINAi, YA.G. and STEPIN, A.M.: 'Theory of
dynamical systems and general transformation groups with
invariant measure', J. Soviet Math. 7, no. 6 (1977), 974-1065.
(Itogi Nauk. i Tekhn. Mat. Analiz 13 (1975), 129-262)
D. V Anosov
References
[AI) KRENGEL, U.: Ergodic theorems, W. de Gruyter, 1985.
[A2) PETERSON, K.: Ergodic theory, Cambridge Univ. Press, 1983.
AMS 1980 Subject Classification: 28DXX
BIRKHOFF - WITf THEOREM,
Poincare- Birkhoff- Witt theorem - A theorem about
the representability of Lie algebras in associative alge-
bras. Let G be a Lie algebra over a field k, let U(G) be
its universal enveloping algebra, and let B = {b i : i El}
be a basis of the algebra G which is totally ordered in
some way. All the possible finite products b", ... b".,
where al ~ . . . ~a" then form a basis of the algebra
U(G), and it thus follows that the canonical homomor-
phism G-"?U(G) is a monomorphism.
It is possible to construct a Lie algebra L(R) for any
multiplication in R with the commutator operation
[xy] = xy-yx.
The Birkhoff - Witt theorem is sometimes formulated
as follows: For any Lie algebra G over any field k there
exists an associative algebra R over this field such that
Gis isomorphically imbeddable in L(R).
The first variant of this theorem was obtained by H.
Poincare [1]; the theorem was subsequently completely
demonstrated by E. Witt [2] and G.D. Birkhoff [3].
The theorem remains valid if k is a principal ideal
domain [4], in particular for Lie rings without opera-
tors, i.e. over Z, but in the general case of Lie algebras
over an arbitrary domain of operators the theorem is
not valid [5].
References
[I) POINCARE, H.: 'Sur les groupes continus', Trans. Cambro Phi-
los. Soc. 18 (1900), 220-225.
(2) WITT, E.: 'Treue Darstellung Liescher Ringe', J. Reine Angew.
Math. 177 (1937), 152-160.
(3) BIRKHOFF, G.D.: 'Representability of Lie algebras and Lie
groups by matrices', Ann. of Math. (2) 38, no. 2 (1937), 526-
532.
[4) LAZARD, M.: 'Sur les algebres enveloppantes universelles de
certain algebres de Lie', c.R. Acad. Sci. Paris Ser. I Math. 234
(1952), 788-791.
(5) SHIRSHOY, A.I.: 'On representations of Lie rings in associative
rings', Uspekhi Mat. Nauk 8, no. 5 (1953), 173175 (in Rus-
sian).
(6) COHN, P.M.: Universal algebra, Reidel, 1981.
(7) KUROSH, A.G.: Lectures on general algebra, Chelsea, 1963
(translated from the Russian).
[8) SERRE, J.-P.: Lie algebras and Lie groups, Benjamin, 1965
(translated from the French).
401
BIRKHOFF - WITI THEOREM
[9] CARTAN, H. and ElLENBERG, S.: Homological algebra, Prince-
ton Univ. Press, 1956.
T.S. Fofanova
AMS 1980 Subject Classification: 17815
BIRTH-AND-DEATH PROCESS - A Markov process
with states 0, 1, ... , in which in a time interval
(t, t +h) transitions from the state n into the states
n + 1 and n -1 occur with probabilities An (t)h +o(h)
and J.'n(t)h +o(h), respectively, and where the probabil-
ity of other transitions is o(h). For a special choice of
the reproduction coefficients An(t) and the death coeffi-
cients J.Ln(t) one can obtain particular cases which pro-
vide a satisfactory description of various real processes:
radioactive transformations, the running of telephone
exchanges, the evolution of biological populations, etc.
The fact that birth-and-death processes are widely used
in applications is due to the simplicity of the equations
for the transition probabilities, which often can be
obtained explicity. For instance, in the case of a Pois-
son process An(t)=A, J.'n(t)=O, the probabilities Pn(t)
(Pn(t) is the probability of being at time t in the state n
if the process has started from state 0) satisfy the sys-
tem:
P'o(t) = -APo(t)
P~(t) = -APn(t)+APn-\(t), n ~ 1,
where
I if n =0,
Pn(O) = {
o otherwise.
The solution of this system is:
Pn(t) = ~e-N,
n.
n =0,1, ....
A more general process is the one in which An(t)=nA,
J.'n(t)=O. This type of process was first studied by G.
Yule (1924) in connection with the mathematical theory
of evolution. A Yule process is a particular case of a
pure birth process which is obtained from the general
birth-and-death process by assuming An (t) = An,
J.'n(t)=O. If An increases very rapidly, then with positive
probability one can pass through all states in a finite
time and then
n =0
The condition
n =0
for a pure birth process is satisfied if and only if the
series ~l IAn diverges. If An(t)=nA+v, Pn(t)=np., then
the birth-and-death process is a branching process with
immigration in which the state n designates the number
of particles. Each particle dies during time (t, t + h)
with probability J.Lh + o(h), splits into two particles with
probability Ah + o(h) and, moreover, one particle immi-
grates from the outside with probability vh +o(h). If
402
v =0 then one obtains the simplest branching process
without immigration. If A=O and v>O, then this type of
process with immigration can be applied to describe the
functioning of a telephone exchange with an infinite
number of lines. Here, the state is the number of occu-
pied lines. The reproduction coefficient An(t)=V charac-
terizes the rate of incoming telephone calls, and J.L is the
expected duration of a conversation.
References
[I] FELLER, W.: An introduction to probability theory and its appli-
cations, 1-2, Wiley, 1957-1971.
[2] SEWASfYANOW, B.A. [BA SEVAST'YANOV]:
Verzweigungsprozesse, Akad. Wissenschaft. DDR, 1974
(translated from the Russian).
[3] SAATY, R.L.: Elements of queueing theory with applicatiOns,
McGraw-Hill, 1961. v.p
ChISfl.) k
.Ta ov
Editorial comments.
References
[A1] ATIlREYA, K.B. and NEY, P.E.: Branching processes,
Springer, 1972.
AMS 1980 Subject Classification: 60J80
BISECTOR PLANE of a dihedral angle, bisector - The
plane passing through the edge of the dihedral angle
and bisecting this angle.
AMS 1980 Subject Classification: 51 NXX
BISECTRIX of an angle - The half-line (ray) issuing
from the apex of the angle and bisecting it. In other
words, the bisectrix of an angle is the set of points
located inside the angle and equally distant from both
of its sides. A bisectrix of a triangle is the segment (and
also its length) of the bisectrix of an internal angle of
the triangle from the apex to the point of intersection
with the opposite side. A bisectrix of a triangle divides
a side of the triangle into segments that are propor-
tional to the adjacent sides. The bisectrices of a triangle
intersect at one point, which is the centre of the
inscribed circle in the triangle. The quadruple of points
A, B, K, L, consisting of two apices A, B of the triangle
ABC, the point K of intersection of the bisectrix of the
angle C with AB, and the point L of intersection of the
bisectrix of the external angle C with AB forms a har-
monic quadruple of points.
A.B. Ivanov
AMS 1980 Subject Classification: 51 NXX
BIT - A binary unit of information, numerically
equal to the amount of information obtained during a
trial with two mutually exclusive equally probable alter-
natives (PI =P2 = I /2):
I(p\,P2) = P\lo~P\ -p2Io~p2 = lo~ 2 = 1 'bit',
the logarithms being taken to the base 2. The bit is the
most frequently used unit, but other units of informa-

tion are also employed - the 'Hartley' or the 'nit', the
definitions of which involve decimal and natural loga-
rithms respectively.
A. V. Prokhorov
Editorial comments. The definition as given comes from
information theory. In computer science the term 'bit' usu-
ally refers to the representation of '0' or '1' by a suitable
physical device which can be (exclusively) in two alternative
states, or even to that device itself. Nits or Hartley's are
unknown in the West.
AMS 1980 Subject Classification: 94A17 , 68A05
BrrsADZE EQUATION - The partial differential
equation that can be written in complex form as fol-
lows:
where w(z)=u + iv, z =x +iy, and that can be reduced
to the elliptic system
U xx - u,y - 2vxy = 0,
vxx-vy'y+2uxy = 0,
in the real independent variables x and y. The homo-
geneous Dirichlet problem in a disc C: I z - z 0 I <t:,
where the radius t: is as small as one pleases, for the
Bitsadze equation has an infinite number of linearly
independent solutions [1]. The Dirichlet problem for
the inhomogeneous equation Wzz = f in the disc C is
normally solvable according to Hausdorff, since it is
neither a Fredholm problem nor Noetherian; in a
bounded domain containing a segment of the straight
line y = 0, this problem is not even a Hausdorff prob-
lem, even though the homogeneous problem has only
one zero solution [2].
References
[I] BrrSADZE, A.V.: 'On the uniqueness of the solution of the Diri-
chlet problem for elliptic partial differential operators', Uspekhi
Mat. Nauk 3, no. 6 (1948), 211-212 (in Russian).
[2] BITSADZE, A. V.: Boundary value problems for second-order ellip-
tic equations, North-Holland, 1968 (translated from the Rus-
sian).
[3] MIRANDA, C.: Partial differential equations of elliptic type,
Springer, 1970 (translated from the Italian).
[4] BERS, L., JOHN, F. and ScHECHTER, M.: Partial differential
equations, Interscience, 1964. A.M. Nakhushev
AMS 1980 Subject Classification: 35J15
BIVECfOR - A class [u, v] of ordered pairs u, v of
vectors of an affine space A, starting at a common ori-
gin (considered in a basis of the underlying space). A
bivector is considered to be equal to zero if its consti-
tuent vectors u and v are collinear. A non-zero bivector
generates a unique two-dimensional space in A, its car-
rier. Two bivectors are said to be parallel if their car-
rier planes are parallel. If A has finite dimension n, and
(u 1 , ,un) are the contravariant coordinates of u,
while (v 1 , ,vn) are the contravariant coordinates of
BIVECTOR
v, calculated with respect to some basis e = (el, . . . , en)
of the underlying space of A, then the quantities
are called the PlUcker coordinates of the pair u, v. Two
pairs of vectors are in the same class if their Plucker
coordinates with respect to some basis coincide (they
will then be equal in any basis). The coordinates of the
class are then called the coordinates of the bivector [u, v]
with respect to the basis e. These coordinates are skew-
symmetric with respect to their indices; they contain (D
independent coordinates. Under a transition to another
basis of A, the coordinates of a bivector behave as
coordinates of a twice-contravariant tensor. A bivector
is also called a free bivector. In the presence of a scalar
product in A, a number of metrical concepts of vector
algebra can be extended to bivectors. The measure of a
bivector is the area of the parallelogram formed by the
vectors u, v, -u, -v, the origin of each one being
located in the end of the precedirig one. This only
depends on the class, not on the representatives u,v.
The scalar product of two bivectors is the number
equal to the product of the measures of the factors by
the cosine of the angle between their two carrier planes.
This product is a bilinear form of the coordinates of
the factors, the coefficients of which are defined by the
metric tensor of the space A alone.
If the dimension of A is 3, the bivector [u, v] may be
identified with a vector of A which, in the presence of a
scalar product, is called the vector product of the vec-
tors u, v.
In tensor calculus a bivector is an arbitrary contra-
variant skew-symmetric tensor of valency 2 (i.e. a ten-
sor of type (2,0. Each such tensor may be represented
as a sum of tensors, to which correspond non-zero
bivectors in the above sense with different carrier
planes. They define the sheets of the bivector. The rank
of the skew-symmetric matrix of dimension n X n con-
sisting of the coordinates of a bivector is an even
number 2p, where p is the number of sheets of the
bivector. In a real affine space A this matrix is similar
to the matrix
o 0
with the blocks
See also Exterior product; Poly-vector; PlUcker coor-
dinates.
403
BIVECTOR
References
[1) ScHOUTEN, JA.: Tensor analysis for physicists, Cambridge
Univ. Press, 1951.
L.P. Kuptsov
Editorial comments. Assign to a non-zero bivector (u, v)
the plane it generates, i.e. the corresponding pOint of the
Grassmannian of 2 planes in (the underlying vector space
of) A. Then the Plucker coordinates of this element in the
Grassmann manifold can be identified with the Plucker
coordinates of the bivector.
References
[A 1) CARTAN,E.: Geometry of Riemannian spaces. (With notes
and appendices by R. Hermann), Math. Sci. Press, 1983
(translated from the French).
[A2) GolAB, S.: Tensor calculus, Elsevier Sci. Publ. (translated
from the Polish).
[A3] RAsHEWSKI, P.K. [P.K. RAsHEVSKIl1: Riemannsche Geometrie
und Tensoranalyse, Deutsch. Verlag Wissenschaft., 1959
(translated from the Russian).
AMS 1980 Subject Classification: 15A72, 53A45
BIVECI'OR SPACE - A centro-affme space EN
(where N=n(n-I}/2), which may be assigned to
each point of a space An with an affine connection (in
particular, to a Riemannian space Vn ). Consider all
tensors with even covariant and contravariant valencies
at a point of the space An (or Vn ); the covariant and
contravariant indices are subdivided into different
pairs, for each one of which the tensor is skew-
symmetric. Tensors with these two properties are
called bitensors. If each skew-symmetric pair is regarded
as a collective index, the number of new indices will be
N = n (n - I) / 2. The simplest bitensor is the bivector
va{3 = - vfJa ~ Va, a, f3 = 1, . . . ,n; a = 1, . . . ,N.
If, at a point P of An,
t
AI} = [~~: Ag = 2A~aAgl = Af~Agt,
v
then =Agv b , and the set of bivectors assigned to An
(or Vn ) at a given point defines a vector space of
dimension N such that the components satisfy the con-
ditions
*
I Ag I 0, AgA~ = 8~
i.e. this set defines the centro-affine space EN, called
the bivector space at the given point. In Vn the bivector
space may be metrized with the aid of the metric tensor
deC
gab = ga{ry~ = gcrygf30-ga~g{3y,
after which EN becomes a metric space R N .
Bivector spaces are used in Riemannian geometry
and in the general theory of relativity. The bivector
space RN is constructed at a given point of the space
Vn , and different representations of the curvature ten-
404
sor with components R aPr8 , R~H, R~~ and the second-
valency bitensors with components R ab , Rg, R~ are
associated, respectively. The study of the algebraic
structure of the curvature tensor may then be reduced
to the study of the pencil of quadratic forms Roo - Agoo,
the second one of which is non-degenerate ( I gab I=FO).
The study of elementary divisors of this pair results in
a classification of the spaces Vn . If n = 4 (N = 6) and if
the form g af3 has signature (- - - +), then it can be
shown that only three types of Einstein spaces exist.
A bivector may be assigned to each rotation in Vn ;
this means that in RN there corresponds a vector, which
is convenient for the study of infinitesimal transforma-
tions. Essentially, a bivector space is identical with a
biplanar space [2].
References
[1) PETROV, AZ.: New methods in general relativity theory, Mos-
cow, 1966 (in Russian).
(2) NORDEN, AP.: 'On complex representation of tensors of a
biplanar space', Kazan. Gos. Univ. Uchen. Zap. 114 (1954), 45-
53 (in Russian).
A.2. Petrov
Editorial comments. Consider a bivector as represented
by an ordered pair of vectors (u, v) as in [A1], [A4] or the
article bivector. The Plucker coordinates of (u, v),
pii = u i vi - u i Vi, then constitute a bivector as in the article
- - -i .
above. Let u=Au, v=Av, i.e. u =Aju l and similarly for v.
Then the pii transform as pii=2A~Allpkl. Whence the for-
mulas above. Here the square brackets in A~AII are a nota-
tion signifying taking an alternating average. Thus
Al1 Atl = (AJ,At - A~A)) / 2. If certain indices are to be sin-
gled out, i.e. exempt from this averaging process, this is
indicated by II. Thus
A~ I ki'l = (A~' - Af/:' - Ai//ii - A~' + Af/:i + A~i) / 6, and
gab=2ga[ygl{3l~l' cf. above. This is a notation introduced by
R. Bach [A5]. Cf. also Alternation.
In more modern terms what is described here is the bun-
dle of bivectors over An.
A centro-affine space is an affine space with a dis-
tinguished point, i.e. practically a vector space. It is not a
term which is still greatly used.
References
[A1] CARTAN, E.: Geometry of Riemannian spaces. (With notes
and appendices by R. Hermann), Math. Sci. Press, 1963
(translated from the French).
[A2] GoLAB, S.: Tensor calculus, Elsevier Sci. Publ., 1974
(translated from the Polish).
[A3] RASHEWSKI, P.K. [P.K. RASHEVSKIT]: Riemannsche Geometrie
und Tensoranalyse, Deutsch. Verlag Wissenschaft., 1959
(translated from the Russian).
[A4] SCHOUTEN, J.A: Ricci-calculus, Springer, 1954.
[A5] BACH, R.: 'Zur Weylschen Relativitatstheorie und der Weyl-
schen Erweiterung des Krummungsbegriffs', Math. Zeitschr.
9 (1921),110-135.
AMS 1980 Subject Classification: 15A72, 15A75,
53805, 53820, 53830
BJORLING PROBLEM - A problem in the theory of
minimal surfaces (cf. Minimal surface), which is to find

the minimal surface passing through a given non-closed
analytic curve L, with given tangent planes along L.
Bjorling's problem on minimal surfaces is the analogue
of the Cauchy problem of differential equations. The
problem was posed and solved by E.G. Bjorling [1]. A
solution of this problem always exists, is unique, and is
explicitly expressed by the Schwarz formula for
minimal surfaces. The solution of Bjorling's problem
always permits one to find the minimal surface when-
ever one of its geodesic lines or one of its asymptotic
lines or one of its lines of curvature is known. If the
given curve is planar and is a geodesic on the minimal
surface which is desired to be found, then the plane of
the curve L will be a plane of symmetry of the minimal
surface.
References
[1] BJORLING, E.G.: Arch. Grunert IV (1844),290.
[2] DARBoux, G.: Lerons sur la theorie generale des surfaces et ses
applications geometriques du calcul infinitesimal, 1, Gauthier-
Villars, 1887.
[3] BLASCHKE, W.: Vorlesungen aber DifJerentialgeometrie und
geometrische Grundlagen von Einsteins Relativitiitstheorie, 1,
Springer, 1921.
[4] NITSCHE, J.c.c.: Vorlesungen aber Minimalfliichen, Springer,
1975.
I.Kh. Sabitov
AMS 1980 Subject Classification: 49C10, 53A 10
BLASCHKE PRODUCf, Blaschke function - A regu-
lar analytic function of a complex variable z, defined in
the unit disc K = {z: I z I < I} by the finite or infinite
product
(*)
where n is a non-negative integer, and {ad,
k = 1,2, ... , is a sequence of points ak EK\ {O} such
that the product on the right-hand side of (*) converges
(the convergence condition is necessary only for an
infinite product). The Blaschke product was introduced
by W. Blaschke [1], who proved the following theorem:
A sequence {ad of points akEK\ {O} defines a func-
tion of the type (*) if and only if the series
~/1-1 ak I) is convergent. Each factor of the form
bk(z)
I ak I ak- z
= ------,
ak l-ak z
called the Blaschke factor for ab defines a univalent
conformal mapping of the disc K onto itself, which
takes z =ak to zero, with the normalization
bk ( -ak / I ak 1)= 1. The factors of the form bo(z)=z
may be interpreted as Blaschke factors which
correspond to the zero z = 0 with the normalization
bo(l)= 1. The definition of Blaschke factors and
Blaschke products is readily carried over to a disc of
arbitrary radius, and also to an arbitrary simply-
connected domain, which is conformally equivalent to a
disc.
BLASCHKE PRODUCT
The sequence 0, ... ,0, a 1. a2, . .. (with n zeros),
which is usually written out in non-decreasing order of
I ak I, is the sequence of all zeros of the Blaschke pro-
duct (*) (each zero is written down as many times as its
multiplicity). Thus, Blaschke's theorem describes the
sequences of zeros of all possible Blaschke products.
The product (*) can be regarded as the simplest
bounded holomorphic function in the disc K with a
prescribed sequence of zeros. It converges absolutely
and uniformly inside K, represents a bounded holo-
morphic function I B (z) I < I in K, with angular boun-
dary values of modulus one almost everywhere on aK.
A necessary and sufficient condition for a bounded
holomorphic function f(z) in K, I f(z) I <1, to be a
Blaschke product, is
[In I j(re
2"
lim
r-+l
i8 ) I dO = O.
Blaschke products may be used to give a product
representation of important classes of holomorphic
functions in the unit disc K. Thus, a proof was given
for the following theorem of Blaschke: A sequence {ak}
of points in the disc K is the sequence of all zeros of
some bounded holomorphic function f (z), If (z) I < 1,
in K if and only if the series ~/1-1 ak I) is conver-
gent. Moreover, f (z) can be represented as a product
j(z)=Blz)g(z),
where BIz) is the Blaschke product constructed with
the zeros {ad of the function f (z), while g (z) is a
zero-free holomorphic function in K, I g(z) I ~1, which
can be represented relatively simply using an integral
formula. Apart from the bounded functions, similar
product representations may be constructed for func-
tions of bounded form and for Hardy classes [2] - [4]
(cf. Function of bounded fonn).
The above theory was considerably generalized by
M.M. Dzhrbashyan [5], [6], who constructed infinite
products of a more general nature, which are suitable
for the factorization of much larger classes of mero-
morphic functions. A solution was also found for the
problem of constructing analogues of Blaschke pro-
ducts and Blaschke's theorem for doubly-connected
domains [7] and, in general, finitely-connected [8]
domains. The solution of the problem of constructing
suitable analogues of the Blaschke product for holo-
morphic functions of several complex variables is ren-
dered very difficult by the fact that the zeros of such
functions cannot be isolated.
References
[I] BLASCHKE, W.: 'Eine Erweiterung des Satzes von Vitali tiber
Folgen analytischer Funktionen', Berichte Math..Phys. KI.,
Sachs. Gesell. der Wiss. Leipzig 67 (1915), 194-200.
[2] l'IuwALOW, I.I. [I.I. PIuvALov]: Randeigenschaften anaiytischer
Funktionen, Deutsch. Verlag Wissenschaft., 1956 (translated
from the Russian).
405
BLASCHKE PRODUCT
[3] NEVANLINNA, R.: Ana(ytic fimctions, Springer, 1970 (translated
from the German).
[4] CoUINGWOOD, RF. and LoUWATER, A.J.: The theory of clus
ter sets, Cambridge Univ. Press, 1966.
[5] DzHuASHYAN, M.M.: Integral transforms and representation of
fimctions in a complex domain, Moscow, 1966 (in Russian).
[6] DzHuASHYAN, M.M.: 'The theory of factorization and boun-
dary properties of functions meromorphic in a disc', Russian
Math. Surveys 28, no. 4 (1973), 1-12. (Uspekhi Mat. Nauk 28,
no. 4 (1973), 3-14)
[7] KAs'yANYUK, SA.: 'On functions of classes A and H 8 in an
annulus', Mat. Sb. 42 (84), no. 3 (1957), 301326 (in Russian).
[8] TAMRAZOV, P.M.: 'Conformal-metric theory of doubly con
nected domains and the generalized Blaschke product', Soviet
Math. Dokl. 6, no. 2 (1965), 432-435. (Dokl Akad. Nauk
SSSR 161, no. 2 (1965), 308-311)
P.M. Tamrazov
Editorial comments. Some standard references for
Blaschke products and related subjects are [A1], [A2] and
[A3].
References
[A1] DUREN, P.L.: Theory of HP spaces, Acad. Press, 1970.
[A2] RUDIN, W.: Real and complex analysis, McGraw-Hili, 1974.
[A3] GARNETT, J.B.: Bounded analytic functions, Acad. Press.,
1981.
AMS 1980 Subject Classification: 30D50
BLASCHKE SELECfION TIIEOREM, Blaschke com-
pactness principle - A metric space of convex bodies is
locally compact, i.e. it is possible to select, out of an
infinite set of convex bodies belonging to a given cube,
a sequence which converges to some convex body in
this cube.
The theorem was demonstrated in 1916 by W.
Blaschke [1].
References
[1] BLASCHKE, W.: Kreis und Kugel, Chelsea, reprint, 1949.
A.B. Ivanov
Editorial comments.
References
[A1] KELLY, PJ. and WEISS, M.L.: Geometry and convexity, Wiley,
1979.
AMS 1980 Subject Classification: 52A07
BLASCHKE - WEYL FORMULA - A variant of the
Green formulas for the rotation field y of an infini-
tesimal deformation of a surface with position vector x:
2 fj [xfu-fv-] du dv = :!c(xy dy).
The proof and the idea of applying the Blaschke - Weyl
formula to demonstrate the rigidity of ovaloids is due
to W. Blaschke [1] and to H. Weyl [2]. For other appli-
cations see [3]. The Herglotz formula is an analogue of
the Blaschke- Weyl formula. The formula has been
generalized to the case of infinitesimal deformations of
surfaces in spaces of constant curvature.
406
References
[1] BLASCHKE, W.: 'Ein Beweis fUr die Unverbiegbarkeit geschlos-
sener konvexer Flichen', Gott. Nachr. (1912),607-610.
[2] WEYL, H.: 'Ueber die Starrheit der Eiflachen und konvexer
Polyeder', Sitzungsber. Akad. Wiss. Berlin (1917), 250266.
Also: Gesammelte Abh., Vol. 1, Springer, 1968, pp. 646-662.
[3] EFIMOV, N.V.: 'Qualitative problems in the theory of deforma-
tions of surfaces', Uspekhi Mat. Nauk 3, no. 2 (24) (1948), 47-
158 (in Russian). Translated into German as book.
[4] BLASCHKE, W.: Einflihrung in die Difjerentialgeometrie,
Springer, 1950. M.l. Voitsekhovskii
AMS 1980 Subject Classification: 53A45
BLOCH CONSTANT - An absolute constant, the
existence of which is established by Bloch's theorem. Let
H be the class of all holomorphic functions f (z) in the
disc I z I < 1 such that J' (0) =1. The Riemann surface
of the function f(z) contains on one of its sheets a
largest open disc of radius Bf>O. It was shown by A.
Bloch [1] that
inf {B/ fEH} = B > O.
The most precise known estimate is V3 / 40;;;;;B 0;;;;;0.472
[2]. It follows from Bloch's theorem that the Riemann
surface of an entire function contains single-sheeted
discs of arbitrary radius; this is equivalent to the Picard
theorem.
References
[1] BLOCH, A.: 'Les theoremes de M. Valiron sur les fonctions
entieres et la theorie de l'uniformisation', Ann. Fac. Sci. Univ.
Toulouse (3) 17 (1925), 1-22.
[2] AHI.FORS, L.V. and GRUNSKY, H.: 'Ueber die Blochsche Kon-
stante', Math. Z. 42 (1937), 671-673.
[3] GoLUZIN, G.M.: Geometric theory offunctions of a complex
variable, Amer. Math. Soc., 1969 (translated from the Russian).
E.D. Solomentsev
Editorial comments. For the connection between the
theorems of Bloch and Picard, see e.g. [A 1).
References
[A 1] HEINS, M.: Selected topics in the classical theory of functions
of a complex variable, Holt, Rinehart and Winston, 1962.
AMS 1980 Subject Classification: 30C75
BLOCK DESIGN - A system of subsets of a finite set
which satisfies certain conditions related to the fre-
quency of appearance of pairs of elements of the set in
the subsets of the system. The concept of a block
design arose in the theory of design (planning) of (sta-
tistical) experiments in the nineteen-twenties and
nineteen-thirties, but had been studied as early as the
mid-nineteenth century under the name tactical confi-
gurations. The concept of a block design is a variation
of the concepts of a hypergraph, a net and a complex.
As a rule, in a block design several additional limita-
tions are imposed on the family of subsets. A block
design may be defined by a pair of sets (V, B) where

Bj ~ V, i=I, ... ,b.
The elements of the set V are called the points (treat-
ments) of the block design, or varieties or elements,
while the elements of the set B are called its blocks. The
element ai and the block Bj are incident if ai EBj . The
number 1 Bj 1 of elements incident with Bj is usually
denoted by kj' while the number of blocks incident
with ai is denoted by ri' The number
I {B/ ajEBj , a/EBj} I
is denoted by Ail. The numbers v, b, ri, kj' Ail
(i, 1= 1, ... , v; j = 1, ... ,b) are said to be the parame-
ters of the block design. If ri =r for all i = 1, ... ,v,
kj=k for all j=1, ... ,b and Ai/=A, (V, B) is a bal-
anced incomplete block design (BIB-design) with param-
eters v, b, r, k, A. The meaning of the word 'balanced' is
that the frequencies of appearance of all elements and
pairs of elements in the blocks are respectively equal,
while the word 'incomplete' indicates that, generally
speaking, not all the k-element sets are included in B.
Let exactly m different numbers AI, . . . ,Am be
encountered among the numbers Ail, i, 1= 1, ... ,v, and
let m symmetric association relations be introduced on
the elements of the set V so that the following condi-
tions are satisfied:
1) the set V2 of all pairs of distinct elements of V is
subdivided into m disjoint subsets VI, ... , V~ and, if
(a, a')E VJ, then a and a' are said to be j-associated;
2) 1 {B/ a EBj , a' EBj , (a, a')E Vt} 1 =Ai'
i=1, ... ,m,j=1, ... ,b;
3)1 {a: 3a' (a, a')EVt} 1 =ni' i=1, ... ,m;
4) 1 {an: (an, a)E VJ, (an, a')E Vi, (a, a')E Vt} I =P)k>
and, in view of the symmetry, P)k =P~j,
i, j, k = 1, ... ,m. A block design with the properties 1)
- 4) is said to be a partial balanced incomplete block
design with m types of relations or a PBIB(m)-design.
The rule which specifies the association relation is
called the association deSign. A BIB-design is a
PBIB(1 )-design. An example of a PBIB(2)-design is a
block design which may be represented as the table
1 222333
456456456
789978897
10 11 12 11 12 10 12 10
where any two numbers in the same column are 1-
associated, while any two numbers in different columns
are 2-associated. Here v=12, b=9, r=3, k=4, AI =1,
A2=0, nl =9, n2=2.
pI = IIp)k II = II~ ~II' p2 = IIp7k II
BLOCK DESIGN
To each block design with v elements and b blocks
corresponds an incidence matrix A = I cij II, where
cij = 1 if ai EBj and Cij =0 otherwise, i = 1, ... , v;
j = 1, ... ,b. The theory of block designs considers
problems on the existence and classification as well as
problems involved in the construction of a block design
with given parameters. The parameters of a block
design are related in a certain manner. The following
equations are valid for BIB-designs:
vr = kb, (1)
A(V -1) = r(k -1).
Equation (1) and the relationships
~nj = v-I,
j=1
m
~Ajnj = r(k -1),
i=1
m. { nj if i*j,
~PJ"k
k=1 = nj - 1 if'l -
- j' , l. , .j-- 1
, ... ,m,
njpjk = njP{k = nkPt, i, j, k = 1, ... ,m,
are valid for the parameters of PBIB(m)-designs. The
incidence matrix of a BIB-design satisfies the fundamen-
tal matrix relation
AAT = (r-A)E+V, (2)
where E is the unit matrix of order v and J is the
matrix of order v consisting exclusively of ones. The
existence of a (0, I)-matrix which satisfies condition (2)
is a sufficient condition for the existence of a BIB-
design with the given parameters. The inequality b ~v
follows from (2). A BIB-design for which b = v (i.e. also
r=k) is said to be a symmetric block design or a
(v, k, A)-configuration. The following theorem applies
to symmetric BIB-designs: If there exists a symmetric
BIB-design with parameters v, k, A, then: 1) if v is even,
k - A is a perfect square; 2) if v is odd, the equation
Z2 = (k-A)x2+(_1)(v-I)/2 Ay 2
has a solution in integers x, y, Z, not all of which are
zero. The conditions of this theorem are sufficient for
the existence of a ratiollal matrix A satisfying equation
(2).
A special range of problems involving the existence
of BIB-designs arises in the context of the following
11, problem: Given b blocks, find the conditions necessary
for the completion of these blocks to form a BIB-
design. In their most general form these conditions are
expressed as the requirement of positive definiteness of
some quadratic form Q, as well as the possibility of
representing Q as a sum of squares of linear forms with
= II~ ~II non-negative coefficients.
The following subclasses of BIB-designs have been
407
BLOCK DESIGN
most extensively studied: Steiner systems (BIB-designs
=
with X 1), in particular Steiner triple systems (k 3); = References
Hadamard configurations (v =b =4t -1, r =k =2t -1,
X=t -1, t ~2), the incidence matrix of which is
obtained from a Hadamard matrix; affine finite
geometries and projective finite geometries [1]. In the
class of PBIB-designs those most extensively studied
are PBIB(2)-designs, which may be subdivided accord-
ing to their association design into group-divisible block
designs, triangular block designs, Latin square block
designs, cyclic block deSigns, etc. [3].
The methods for constructing block designs are usu-
ally classified as direct and recursive. The latter methods
make it possible to use designs with smaller parameter
values to construct designs with larger parameter
values. The direct methods mostly utilize the properties
of finite fields or some geometric properties.
Block designs are used in the design of experiments,
the theory of games, graph theory and in the construc-
tion of error-correcting codes.
References
[I] RySER, R.J.: Combinatorial mathematics, Wiley, 1963.
[2] HAll, M.: Combinatorial theory, Blaisdell, Waltham (MA),
1967.
[3] SHIROKOVA, S.A.: 'Block designs', Russian Math. Survrys 23,
no. 5 (1968), 47-94. (Uspekhi Mat. Nauk 23, no. 5 (1968), 51-
98) V.E. Tarakanov
Editorial comments. The theorem giving conditions on
the numbers v, k, A in order that a symmetric BIB-design
with these numbers as parameters exists, is called the
Bauck- Ryser- Chawla theorem.
AMS 1980 Subject Classification: 05B05
BLOCK-DIAGONAL OPERATOR with respect to a
given orthogonal decomposition H = ~k;;'1 (f) H k of a
Hilbert space H - A linear operator A on H which
leaves each of the subspaces H k , k~l, invariant. The
spectrum of A is the closure of the union of the spectra
of the 'blocks' A IH. =Ak> k~l, II A I =SUPk;;.1 I Ak II
A block-diagonal operator in the broad sense of the
word is an operator A of multiplication by a function X
in the direct integral of Hilbert spaces
H = L$H(t)dp.(t), (Aj)(t) = A(t)f(t), t EM.
Here X(t) is a linear operator acting on the space H(t).
Each operator which commutes with a normal operator
is a block-diagonal operator with respect to the spectral
decomposition of this operator. See also Diagonal
operator.
References
[1] NAlMARK., M.A.: Normed rings, Reidel, 1984 (translated from
the Russian).
N.K. Nikol'skii
B.S. Pavlov
408
Editorial comments.
[A1] lIALMos, P.R.: A Hilbert space problem book, Springer,
1982.
AMS 1980 Subject Classification: 47BXX
BLOITO GAMES - A class of two-person zero-sum
games (cf. Two-person zero-sum game) in normal form
in which the pure strategies (cf. Strategy (in game
theory) of the players are distributions of limited
resources (indivisible or divisible) over several objects,
and the gain, or payoff, is equal to the sum of the gains
on the individual objects. So named after the fictional
colonel Blotto who was supposed to have participated
in one of the first games of this type.
References
[1] KARLIN, S.: Mathematical methods and theory in games, pro-
gramming and economics, Addison-Wesley, 1959.
I.N. Vrublevskaya
AMS 1980 Subject Classification: 90D05
BOCHNER ALMOST-PERIODIC FUNCfIONS -
Functions equivalent to Bohr almost-periodic functions;
defined by S. Bochner [1]. A function f (x) which is
continuous in the interval (- 00, (0) is said to be a
Bochner almost-periodic function if the family of func-
tions {f (x + h ): - 00 < h < oo} is compact in the sense
of uniform convergence on (- 00, (0), i.e. if it is possi-
ble to select from each infinite sequence f (x + hk ),
k = 1, 2, ... , a subsequence which converges uniformly
to f (x) on (- 00, (0). Bochner's definition is exten-
sively employed in the theory of almost-periodic func-
tions; in particular, it serves as the starting point in
abstract generalizations of the concept of almost-
periodicity.
References
[I] BOCHNER, S.: 'Beitrage zur Theorie der fastperiodischen Funk-
tionen I, Funktionen einer Variablen', Math. Ann. 96 (1927),
119-147.
[2] LEVITAN, B.M.: Almost-periodicfunctions, Moscow, 1953 (in
Russian).
E.A. Bredikhina
Editorial comments.
References
[A1] MAAK, W.: Fastperiodische Funktionen, Springer, 1967.
[A2] AMERIO, L. and PROUSE, G.: Almost-periodic functions and
functional equations, v. Nostrand, 1971.
AMS 1980 Subject Classification: 42A75, 43A60
BOCHNER INTEGRAL - An integral of a function
with values in a Banach space with respect to a scalar
measure. It belongs to the so-called strong integrals
(cf. Strong integral).
Let F(X; E, 18, JL) be the vector space of functions
x (f), fEE, with values in a Banach space X, given on a
 BOCHNER - MARTINELLI REPRESENTATION FORMULA

space (E, ~,p.) with a countably-additive scalar meas- and

ure p. on a a-algebra ~ of subsets of E. A function xn(t)dpo ~ x(t)dp..
Xo EF is called simple if
4) The space .fl? is complete with respect to the norm
(d. Convergence in nonn)
_ {Xi, tEBiE'B, p.(Bi)<oo, Bi nBj= 0,
xo(t) - N'=I= '. . .-1 N II x-y II = III x(t)-y(t) II dp..
0, tEE \ UBi' 1 J, I, } - , . . . , , B
i=l
5) If T is a closed linear operator from a Banach
A function x EF is called strongly measurable if there space X into a Banach space Y and if
exists a sequence {xn} of simple functions with X E Ie(X; E, '13, po), Tx E Ie(Y; E, '13, po),
I x - Xn II ~O almost-everywhere with respect to the then
measure p. on E. In such a case the scalar function TX(t)dpo = TX(t)dp..
II x II is ~-measurable. For the simple function Xo
def N If T is bounded, the condition
ixo(t)dpo = i~Xip.(B;).
Tx E Ie(Y;E, 'B,po)
A function x is said to be Bochner integrable if it is is automatically fulfilled, [3] - [5].
strongly measurable and if for some approximating The Bochner integral was introduced by S. Bochner
sequence {xn} of simple functions [1]. Equivalent definitions were given by T. Hildebrandt
lim I II X (t) - xn(t) I dpo = O. [2] and N. Dunford (the Do-integral).
n-+oo E
References
The Bochner integral of such a function over a set B E ~ [I] BOCHNER, S.: 'Integration von Funktionen, deren Werte die
Elemente eines Vektorraumes sind', Fund Math. 20 (1933),
is
IX (t) dpo = lim
B n-+oo E
I XB(t)Xn(t) dpo, 262-276.
[2] HILDEBRANDT, T.H.: 'Integration in abstract spaces', Bull.
Amer. Math. Soc. 59 (1953), 111-139.
where XB is the characteristic function of B, and the [3] YOSIDA, K.: Functional analysis, Springer, 1980.
limit is understood in the sense of strong convergence [4] HILLE, E. and PHILLIPS, R.: Functional analySiS and semi-
in X. This limit exists, and is independent of the choice groups, Arner. Math. Soc., 1957.
[5] DUNFORD, N. and ScHWARTZ, J.T.: Linear operators. General
of the approximation sequence of simple functions. theory, I, Interscience, 1958.
Criterion for Bochner integrability: For a strongly- v.l. Sobolev
measurable function x EF to be Bochner integrable it is Editorial comments. A simple function is also called a
necessary and sufficient for the norm of this function to step function. A good recent textbook on integrals with
be integrable, i.e. values in a Banach space is [A1]; [A4] is specifically about
the Bochner integral.
III x(t) II dpo < 00.
B References
[A1] DIESTEL, J. and UHL, J.J.: Vector measures, Math. Surveys,
The set of Bochner-integrable functions forms a vec- 15, Amer. Math. Soc., 1977.
tor subspace .fl? of F, and the Bochner integral is a [A2] ZAANEN, A.C.: Integration, North-Holland, 1967.
linear operator on this subspace. [A3] BOURBAKI, N.: Elements of mathematics. Integration,
Addison-Wesley, 1975, Chapt.6; 7; 8 (translated from the
Properties of Bochner integrals: FrenCh).
1) [A4] MIKUSrNSKI, J.: The Bochner integral, Academic Press, 1978.

AMS 1980 Subject Classification: 28805, 46G10

2) A Bochner integral is a count ably-additive p.- BOCHNER. - MARTINELLI REPRESENTATION

absolutely continuous set-function on the a-algebra ~, FORMULA, Bochner- Martinelli representation,
1.e. 00 Bochner- Martinelli formula - An integral representa-
U~IB,X(t)dpo = i~jX(t)dpo, tion for hoIomorphic functions, which is defined as fol-
lows [1], [2]. Let the function f be holomorphic in a
if B j E~, B j n B = 0, i=/=}, p.(B
j j) < 00, and domain D ccn with piecewise-smooth boundary aD,
II kx(t)dp. II~O if p.(B)~O, uniformly for BE'Jj. and let f be continuous in its closure D. Then the
3) If Xn EF, Xn~X almost-everywhere with respect to expresslOn
(n -l)! ( ren x (1)
the measure p. on B E'Jj, if I Xn II.;;;;f almost-everywhere (2'TTi)" iD I r- z Iln
with respect to p. on B, and if fJ"(t) dp.< 00, then n _ _ _
x ~ (rj - Z)drl /\ drl /\ ... /\ [drjl/\ drj /\ ... /\ drn =
X E Ie j=l

409
BOCHNER -MARTINELLI REPRESENTATION FORMULA

= {fez), if ZE~ boundary values of holomorphic functions of several

0, if Z ~D, complex variables. It was introduced by S. Bochner [1]
and by E. Martinelli [2].
here [dIj] means that the term dIj is to be omitted, is References
called the Bochner- Martinelli representation. For [1] BocHNER, S.: 'Analytic and meromorphic continuation by
n = 1 this representation is identical with the Cauchy means of Green's formula', Ann. of Matk (2) 44, no. 4 (1943),
integral formula (cf. Cauchy integral), but for n > 1 its 652-673.
[2] MARTINELLI, E.: Rend Accad Ital. 9 (1938),269-283.
kernel is not holomorphic in z, which is the reason for [3] VLADIMIROV, V.S.: Methods of the theory offunctions of several
the limited applicability of the Bochner- Martinelli complex variables, M.l.T., 1966 (translated from the Russian).
representation in the theory of functions of several E.M. Chirka
complex variables. The kernel of the Editorial comments. SokhotskIT's formula is also known
Bochner- Martinelli representation is the ~-differential as Plemelj's formula. The Cauchy-Green formula is also
form of bidegree (n, n -1): known as the Cauchy- Pompeiu formula.
w Z - (n -I)! 1 x The Bochner-Martinelli kernel is a special case of a
(~, ) - (2"'i)" I ~ - Z 12n Cauchy- Fantappie kernel. The integral representation (2)
solves the a-equation:
X ~(Ij-Zj)dII /\ d~1 /\ ... /\ [dIjl/\d~j /\ ... /\ dIn /\ d~n'
j=1 af =v with av = 0,
which is defined in en, has a singular point at ~=z, for compactly-supported (0, 1)-forms v, by substituting v for
and is a-closed (i.e. aw = 0) outside the singular point. d f in the right-hand side of (2) and omitting the integral over
If n > 1, the form w is equal to aw' (~, z), where the boundary. When v is not compactly supported, the

(~, z) = - (~2-i! .I ~ _ 112n -2~ [IIdIk /\ d~k]

boundary integral causes difficulties. These can be solved
w' for strictly pseudo-convex domains and the
"'/ Z J=I k~J Bochner-Martinelli kernel then occurs in an explicit solu-
is a form of bidegree (n -I, n -1), the coefficient of tion operator for the a-equation.
which is a fundamental solution of the Laplace equa- References
tion; here [A 1]KRANTZ, S.: Function theory of several complex variables,
Wiley, 1982.
d'" - d'"
d</> = ~dZk ~ and d</>
uZk
= ~dZk-;-:r--.
uZk
[A2] HENKIN, G.M. and LEITERER, J.: Theory of functions on com-
plex manifolds, Birkhauser, 1984.
The following integral representation, which general- [A3] RANGE, R.M.: Holomorphic functions and integral represen-
tations in several complex variables, Springer, 1986.
izes formula (1), is the analogue of the Cauchy - Green
formula (cf. Cauchy integral): If the function f is con- AMS 1980 Subject Classification: 32A25
tinuously differentiable in the closure of a domain
D ccn with piecewise-smooth boundary aD, then, for BOGOLYUBOV CHAIN OF EQUATIONS, BBGKY-
any point zeD, equations, Bogolyubov - Born - Green - Kirkwood - Yvon
fez) = !J<nwc~, z)- ,ff<n/\wct z). (2)
equations - A chain of equations (hierarchy) for the
one-particle, two-particle, etc., distribution functions of
The function a classical statistical system. The functions are defined
fez) = ff<nw(~, z), as follows:
where f is a smooth hypersurface in R2n =en and f is a F,(t, rl> ... ,r" PI> ... ,pJ = (1)
function on f which is Lebesgue-integrable, is said to V' jwNdrs+1 ... drNdp,+1 ... dpN,
be an integral of Bochner- Martinelli type. As for
Cauchy-type integrals, SokhotskiI's formula, with the where s = I, 2, ... , V is the volume of the system and
usual restrictions on f and f, is applicable to WN is the N-particle normalized distribution function
Bochner - Martinelli-type integrals. A which satisfies the Liouville equation:
Bochner- Martinelli-type integral is a complex function OWN
at === {H, WN}, (2)
which is harmonic everywhere outside f; in the general
case this function is holomorphic only for n = I. If where the brackets are the Poisson brackets and H is
f=aD, then if n~l, the conditi(~.nJ(z) 0 outside Dis the Hamiltonian of the system. In the thermodynami-
equivalent to the holomorphy of finD. cal limit V~oo, V / N=v=const, the chain of equa-
The Bochner- Martinelli representation is employed tions has the form of an equation for Fs with a specific
to demonstrate other integral representations (e.g. the 'link' with the function Fs + 1 of a higher rank:
Bergman - Weil representation), in holomorphic con- of,
tinuation from the boundary, and also in the theory of at - {H" F,} === (3)

410

= -1 j {
V
~
I <;;;,.;s
<P( I r, -rs+1 I),Fs+I}drs+1 dps+h
where <P( I ri - rs + I I) is the potential of the interaction
between the i-th and the (s + l)-th particles, while Hs is
the Hamiltonian of the s-particle system.
Under conditions of thermodynamic equilibrium
when the momentum distribution of each particle is the
MaxweU distribution, one considers the s-particle coor-
dinate distribution functions of particles, defined by
relations of type (1) in terms of an N-particle function
WN(rh,rN) = ~exp[-~I},
where HI = H - H 0, and H 0 is the sum of the kinetic
energies of the particles of the system, while the. co~i
guration integral Q is defined by the normalizatIOn
condition (4). The chain of equations for these func-
tions has the form
-aFs
+ -I-aus
F+
aXI 0 aXI S
Ilja<p(lrl-rs+II)F d
+-0 -v a
XI
s+1 rs+1 = 0,
where Us is the potential energy of the interaction
between s particles of the system.
The distribution functions, in particular FI and F 2 ,
may be used to express all the specific characteristics of
statistical systems. The principal difficulties involved in
the study of the functions (3) or (5) are caused by the
problems of closing the system (decoupling the chain of
equations) and solving the closed system under special
boundary conditions for the functions Fs. This study is
essential for physical systems of various kinds and is
most advanced for the case of short-range interaction,
when rb / v 1, where ro is the effective radius of
interaction between the particles, and for the case of
long-range interaction, when rb / v 1, in particular
for systems in which electrostatic (Coulomb) forces are
operative. In the time-dependent theory this leads
directly to the kinetic Boltzmann equation for the
single-particle function F I or to the V1asov kinetic
equation, and in the equilibrium theory to the virial
decomposition for the thermodynamic potential or to
specific Coulomb corrections.
In considering quantum-statistical systems, the
BBGKY-hierarchy is constructed for s-particle statisti-
cal quantum operators <x I, . . . ,Xs I Fs Ix~ , ... ,x~ >
which are traces over the variables s + 1, ... ,N of the
particles of the general N-particle operator - the den-
sity matrix. The form of these equations is similar to
that of (3), in which the classical Poisson brackets have
been replaced by quantum brackets.
References
[I] BOGOLYUBov, N.N.: 'Problems of the dynamical theory in sta-
BOGOL YUBOV INEQUALITY
tistical physics', in J. de Boer and G.E. Uhlenbeck (eds.): Stu-
dies in statistical mechanics, Vol. I, North-Holland, 1961.
[2] UHLENBECK, G.E. and FORD, G.V.: Lectures in statistical
mechanics, Amer. Math. Soc., 1963. lA K ko
. . vasm v
AMS 1980 Subject Classification: 82A05
BOGOLYUBOV INEQUALITY in statistical mechan-
ics - 1) Bogolyubov's inequality for the free-energy func-
tional is an inequality that gives rise to a variational
principle of statistical mechanics. The following ine-
quality is valid for any Hermitian operators U I and
(4) U2 :
~<UI-U2>UI o;;;;f[Utl-/[U2] 0;;;; (*)
1
0;;;; N <UI - U2>u"
where
flU] =- ~ In Tre- u /8.
(5)
This expression has the meaning of the free-energy d~n
sity for a system with Hamiltonian U; the extenSIve
parameter N is the number of particles or the volume,
depending on the system; E> is the absolute temperature
in energy units, and
_ Tr( ... e- u /8)
< >u= Tre-U /8
denotes the thermodynamic average with respect to the
Hamiltonian U.
The Bogolyubov inequality (*) is used to obtain the
exact thermodynamic limit solutions of model problems
in statistical quantum physics [1], [2], in studies using
the method of molecular fields [3], in proving the
existence of the thermodynamic limit, and also in order
to obtain physically important estimates for the free
energies of various multi-particle systems [4]. There
exist generalizations of the Bogolyubov inequality (*) to
the case of von Neumann algebras with a 'trace' [5] and
general von Neumann algebras [6].
References
[1) BOGOLYUBOV, JR., N.N.: 'On model dynamical systems in sta-
tistical mechanics', Physica 32 (1966), 933-944.
[2) BOGOLYUBOV, JR., N.N.: A method for studying model Hamil-
tonians, Pergamon Press, 1972 (translated from the Russian).
[3] TYABLlKOV, S.V.: Methods of the quantum theory of magnetism,
Plenum Press 1967 (translated from the Russian).
[4) KUDRIN, L.P:: Statistical plasma phySiCS, Moscow, 1974 (in
Russian).
[5) RUSKAI, M.B.: 'Inequalities for traces on von Neumann alge-
bras', Comm. Math. Phys. 26 (1972),280-289. .
[6) AIwa, H.: 'Golden - Thompson and Peierls - Bogolyubov me-
qualities for a general von Neumann algebra', Comm. Math.
Phys. 34 (1973), 167-178.
Editorial comments.
References .
[A1] RUELLE, D.: Statistical mechanics: rigorous results, Benjamin,
1974.
411
BOGOLYUBOV INEQUALITY
2) Bogolyubov's inequality for Green functions and
correlation functions. The following inequality holds for
time-temperature commutator Green functions (cf.
Green function in statistical mechanics). If one defines
I f(z), z=(zJ, ... ,zn)=x+ry, be holomorphic in an
I open set ~ = {z: I z I <7J, Y EC}, where C is an open
8
A;B = 2~ <AB;s>Hds,
where Bis is the operator B in the Heisenberg represen-
tation for imaginary time t = is, < ... > denotes the
thermodynamic average for Hamiltonian H and () is the
absolute temperature, then
I A;B 12 os;;;
os;;; I A;A+ 1'1 B+;B I
Substituting A =iQ=[Q, H]_ one gets
I B+;B I os;;; I <[Q,Bl_>HI2
277' I <[Q, [Q+, Hl-l- >H
where [,]_ is the commutator. One may also write
down the inequality
<BB++B+B> ~ 28 I <[Q, Bl->H 12
...- I <[Q,[Q+,Hl-l->HI
which follows from (2). Owing to the general nature of
the inequalities (2) and (3) they are extensively
employed in studies of various physical systems.
An improved estimate of the correlation function
<BkBt +Bt Bk>H
is attained in (3) by selecting as the operator Qk = Q (k)
some 'quasi-integral' of motion which, for k=O, com-
mutes with the Hamiltonian H. The commutator in the
numerator of the right-hand side of (3) describes the
transformation properties of Bk under the infinitesimal
transformations of the continuous symmetry group gen-
erated by the operator Qk=O' The inequalities (2) and
(3) are effectively employed in the study of systems
with a spontaneous symmetry breaking: In such a case
thermodynamic averages must be considered in the
framework of the method of quasi-averages (cf. Quasi-
averages, method of).
Similar inequalities are applicable to Green functions
in classical statistical mechanics, where the respective
commutators 'become' Poisson brackets.
Bogolyubov's inequalities made it possible to estab-
lish several relations concerning model systems of sta-
tistical physics, to study the problem of ordering in
infinite systems, etc. For references see Bogolyubov
theorem.
A.M. Kurbatov
AMS 1980 Subject Classification: 82A 15
BOGOLYUBOV THEOREM - I) Bogolyubov's edge-
of-the-wedge theorem is a generalization of the principle
of analytic continuation, in particular to the case of
412
several complex variables. It was obtained in 1956 by
N.N. Bogolyubov in the justification of the dispersion
relations in quantum field theory ([1], Appendix A).
The modem formulation is as follows. Let a function
cone in Rn with apex at zero such that C C)= 0 ,n (-
let the open set 0 eRn be contained in the ball
I x I <7J and suppose that for any test function <p(x)
from D( 0) the limit
lim jf(x +iy)cp(x)dx
y-+O
(1) yeC
exists, independent from the way in which y ~O, Y E C;
then f (z) caI'!. be analytically continued into the
domain r;; u
0:
(2) o = U{z: I z-~ I <8~o(~1,
I'
where 0 is a complex neighbourhood of the set 0,
O<()< I is a constant which depends only on the cone
C, and ~om is the distance from the point ~ to the
(3)
boundary of O. Bogolyubov's edge-of-the-wedge
,
theorem also remains valid if 7J = 00. In such a case,
and under certain assumptions regarding the growth of
the function f (z), one obtains the original formulation
of Bogolyubov (cf. [1]; the light cone in R4 plays the
role of the cone). There exist various proofs and gen-
eralizations of this theorem [2], [5]. In particular, one
can mention the generalizations to hyperfunctions [4]
and to holomorphic cocYcles [3].
Bogolyubov's edge-of-the-wedge theorem is exten-
sively employed in axiomatic quantum field theory, in
the theory of partial differential equations and in the
theory of boundary values of holomorphic functions
(especially of functions of several complex variables). A
useful completion of the theorem is the C-convex hull
theorem [2]. Let, under the conditions of Bogolyubov's
theorem, C=C+ U C-, C- = -C+, where C+ IS a
convex sharp cone; then
Bc(O) C ReH(T;i U i)),
where H(D) is the envelope of holomorphy (cf. Holo-
morphic envelope) of D, ReD is the real section of the
domain D, and Be(O) is the C-convex hull of the set
0, i.e. the smallest open set containing 0 and having
the following property: If the points x' and x" of
Be(O) can be connected by a C-similar curve that is
totally contained in Be(O), then all C-similar curves
homotopic to it are located in Be( 0).
References
[I] BOGOLYUBOV, N.N. and SHIRKOV, D.V.: Introduction to the
theory of quantized fields, Interscience, 1959 (translated from
the Russian).
[2] VLADIMIROV, V.S.: Methods of the theory offunctions of sl!\Ieral
complex variables, M.LT., 1966 (translated from the Russian).

(3) MARTINEAu, A.: 'Distributions et valeurs au bord des fonc-
tions holomorphes', in Theory of distributions. Proc. internat.
summer inst., Lisbon, 1964, pp. 193-326.
(4) KOMATSU, H. (ED.): Hyperfonctions and pseudodifferentiaJ equa-
tions. Proc. Con! Katata, 1971, Springer, 1973.
(5) RUDIN, W.: Lectures on the edge-ofthe-wedge theorem, Amer.
Math. Soc., 1971.
V.S. Vladimirov
Editorial comments. Additional generalizations of the
edge-of-the-wedge theorem include the case where
Cn (- C)= 0, see [5], and the case where the 'edge' of
the 'wedge' is a totally-real hypersurface in en
instead of an
open set in {lmz=O} in en,
see [A1]. For cones
C= C+ U C- , where C+ is convex, a C-similar curve in 0
is defined as a differentiable curve C(t): [0, 1]~O such that
C '(t)EG.
References
[A1] BEDFORD, E.: 'Holomorphic continuation at a totally real
edge', Math. Ann. 230 (1977), 213-225.
2) Bogolyubov's theorem on singularities of type 1/ q2.
A theorem in statistical mechanics on the asymptotic
behaviour of Green functions (cf. Green function) in
the limit of small momenta (q~O) for Bose and Fermi
systems with a gauge-invariant interaction potential.
Presented in 1961 by N.N. Bogo1yubov [1].
The following inequality is valid for systems of
several interacting particles in the case of a degenerate
statistical equilibrium state for the two-time tempera-
ture commutator Green functions (cl. Green function in
statistical mechanics) in the energy representation:
I a q , at E=O I ~ constq-2, (1)
a at
where q , are the annihilation and creation opera-
tors of a particle with momentum q.
The singularities in the Green functions specified in
Bogo1yubov's theorem which appear when q~O
correspond to elementary excitations in the physical
system under study. Bogo1yubov's theorem also predicts
the asymptotic behaviour for small momenta of
macroscopic properties of the system which are con-
nected with Green functions by familiar theorems.
Thus, according to (1), in the case of superfluid Bose
or Fermi systems, the particle density distribution w(q)
tends to infinity as q~O at a rate which is not slower
than 1 / ri. In such a case the degeneracy of the state
of statistical equilibrium is due to the conservation law
of the total number of particles, i.e. to the invariance of
the Hamiltonian of the system with respect to gauge
transformations. However, similar singularities appear
in the respective Green functions, and hence also in the
correlation functions characterizing systems with other
kinds of degeneracy caused by certain additional con-
servation laws, i.e. by the invariance of the Hamil-
tonian of the system with respect to certain transforma-
tions. Bogolyubov's theorem has a number of non-
trivial physical consequences, involving problems of
long-range ordering in systems of many interacting par-
BOHL ALMOST-PERIODIC FUNCTIONS
ticles in which a spontaneous symmetry breaking mani-
fests itself in various ways: the Heisenberg model with
ferromagnetic, anti-ferromagnetic and ferrimagnetic
ordering, systems of superfluid and superconducting
types and systems with crystalline ordering.
The appearance of singularities in the Green func-
tions as q~O is connected with the presence of a
branch of collective excitations in the energy spectrum
of the system that corresponds with spontaneous sym-
metry breaking under certain restrictions on the
interaction potential.
The nature of the energy spectrum of elementary
excitations may be studied with the aid of the mass
operator inequality constructed for Green functions of
type (1)_ In the case of Bose systems, for a finite tem-
perature (8=kB T>0), this inequality has the form:
I ~1I (0, q) - ~12 (0, q) I 0;;;; constq2. (2)
If q = 0, formula (2) yields a generalization of the so-
called Hugenholtz - Pines formula to finite tempera-
tures_ If one assumes that the mass operator is regular
in a neighbourhood of the point (E =0, q= 0), then one
can use (2) to prove the absence of a gap in the
(phonon-type) excitation energy spectrum.
In the case of zero temperatures (8=0) inequality (1)
establishes a connection between the density of the
continuous distribution of the particle momenta and
the minimum energy of an excited state.
Relations of type (1) should also be valid in quantum
field theory, in which a spontaneous symmetry break.-
ing (at a transition between two ground states) results
[4], [5] in an infinite number of particles of zero mass
(Goldstone's theorem), which are interpreted as singular-
ities for small momenta in the quantum field Green
functions_ Bogolyubov's theorem has been applied to a
relativistic quantum-field model with a spontaneous
symmetry breaking [6].
References
(1) BOOOLYUBOV, N.N.: Selected woria, Vol. 3, Kiev, 1970 (in
Russian).
(2) SADOVNIKOV, B.1. and FEDYANIN, V.K.: 'N.N. Bogoljubov's
inequalities in systems of many interacting particles with bro-
ken symmetry', Teoret. Mat. Fiz. 16, no. 3 (1973), 368-393 (in
Russian). English abstract.
(3) BoooLYUBOV, JR., N.N. and SADOVNIKOV, B.I.: Some ques-
tions in statistical mechanics, Moscow, 1975 (in Russian).
(4) GoWSIONE, J.: 'Field theories with 'superconductor' solu-
tions', Nuovo Cimento (10) 19 (1961), 154-164. Italian abstract.
(5) GoWSIONE, J., SALAM, A. and WEINBERG, S.: 'Broken sym-
metries', Phys. Rev. (2) 127 (1962), 965-970.
(6) KAzANSKIi, A.K.: Teoret. Mat. Fiz. n, no. 3 (1975), 418-421.
A.M. Kurbatov
AMS 1980 Subject Classification: 32015, 46F20,
82A15
BOllI.. ALMOST-PERIODIC FUNCfIONS - A class
of functions the typical property of which is that they
413
BOHL ALMOST-PERIODIC FUNCTIONS
can be uniformly approximated on the whole real axis
by generalized trigonometric polynomials of the form
where n 1, . . . ,nkare arbitrary integers, while
al> ... ,ak are given real numbers. This class of func-
tions contains the class of continuous 2'IT-periodic func-
tions and is contained in the class of Bohr aImost-
periodic functions. P. Bohl specified several necessary
and sufficient conditions for a function to be almost-
periodic. In particular, any function of the type
f(x) = fl(X)+ ... +Jk(x),
where each one of the functions f1 (x), ... ,jk(x) is
continuous and periodic (with possibly different
periods), is a Bohl almost-periodic function.
References
[1] BOHL, P.: Uber die Darstellung von Funktionen einer Variabeln
durch trigonometrische Reihen mit mehreren einer Variabeln pro
portionalen Argumenten, Dorpat, 1893. Thesis.
[2] BOHL, P.: 'Ueber eine Differentialgleichung der
StOrungstheorie', J. Reine Angew. Math. 131 (1906), 268-321.
[3] LEVITAN, B.M.: Almost-periodic junctions, Moscow, 1953 (in
Russian).
E.A. Bredikhina
Editorial comments. A very well-known reference for
this kind of topic is [A 1].
References
[A1] BOHR, H.: Almost periodiC functions, Chelsea, 1947
(translated from the German).
AMS 1980 Subject Classification: 42A75, 43A60
BOHR ALMOST-PERIODIC FUNcnONS, uniform
almost-periodic functions - The class U-a.-p. of almost-
periodic functions. The first definition, which was given
by H. Bohr [1], is based on a generalization of the con-
cept of a period: A continuous function f(x) on the
interval (- 00,00) is a Bohr almost-periodic function if
for any (>0 there exists a relatively-dense set of (-
almost-periods of this function (d. Almost-period). In
other words, f (x) is U-almost-periodic (or EO U-a.-p.) if
for any (>0 there exists an L =L() such that in each
interval of length L there exists at least one number T
such that
I f(x+T)-f(x) I <l, -oo<x<oo.
If L(), (~O, is bounded, a Bohr almost-periodic func-
tion f (x) becomes a continuous periodic function.
Bochner's definition (cf. Bochner almost-periodic func-
tions), which is equivalent to Bohr's definition, is also
used in the theory of almost-periodic functions. Func-
tions in the class of U-almost-periodic functions are
bounded and uniformly-continuous on the entire real
axis. The limit f (x) of a uniformly-convergent sequence
of Bohr almost-periodic functions {j,,(x)} belongs to
the class of U-almost-periodic functions; this class is
invariant with respect to arithmetical operations (in
414
particular the Bohr almost-periodic function
f (x) / g(x) is U-almost-periodic, under the condition
inf
-oo<x<oo
I g(x) I > Y > O.
If f (x) is U-almost-periodic and if / (x) is uniformly
continuous on ( - 00, 00), then / (x) is U-almost-
periodic; the indefinite integral F(x)= (f(t)dt is U-
almost-periodic if F(x) is a bounded funcOtion.
References
(1] BOHR, H.: 'Zur Theorie der fastperiodischen Funktionen 1',
Acta Math. 45 (1925), 29-127.
[2] LEVITAN, B.M.: Almost-periodic junctions, Moscow, 1953 (in
Russian).
E.A. Bredikhina
Editorial comments. Bohr's treatise [A1] is a good refer-
ence. An up-to-date reference is [A2].
References
[A 1] BOHR, H.: Almost periodic functions, Chelsea, 1947
(translated from the German).
[A2] CORDUNEANU, C.: Almost periodic functions, Wiley, 1968.
AMS 1980 Subject Classification: 42A75, 43A60
BOHR COMPACTIFICATION - The space X of max-
imal ideals of the algebra of Bohr almost-periodic func-
tions (d. Banach algebra; Bohr almost-periodic func-
tions). The Bohr almost-periodic functions form a
commutative C' -algebra A over the reals R. The alge-
bra A is isometrically isomorphic to the algebra C (X)
of all continuous functions on the compactum X. The
real line R is naturally imbedded in X as an
everywhere-dense subset (however, this imbedding is
not a homeomorphism). The compactum X has the
structure of a connected compact group which is identi-
fied with the group of characters of the real line if the
latter is considered with the discrete topology. This
isomorphism between the algebra of almost-periodic
functions and the algebra of all continuous functions
on the Bohr compactification makes it possible to sim-
plify the proofs of a number of classical theorems. The
concept of a Bohr compactification is also meaningful
for the algebras of almost-periodic functions on other
groups. In the case of the set of conditionally-periodic
functions with n independent fixed periods, the n-
dimensional torus with these periods acts as the Bohr
compactification.
References
[1] LOOMIS, L.H.: An introduction to abstract harmonic analysis, v.
Nostrand, 1953.
E.A. Gorin
Editorial comments. Instead of Bohr compactification
the term Bohr compactum is also used.
There are several other ways to characterize the Bohr
compactification and to prove its existence; see e.g. [A 1].
The most abstract one is the characterization of the Bohr
compactification of a topological group G as the reflection
of G in the category of all compact Hausdorff topological

groups; its existence is then guaranteed by the adjoint func-
tor theorem. This approach is also used for semi-groups
and leads to several other types of compactifications (e.g.
the weakly almost-periodic compactifications). See [A2].
References
[A 1] ALFSEN, E.M. and HOLM, P.: 'A note on compact representa-
tions and almost periodicity in topological groups', Math.
Scand. 10 (1962), 127-136.
[A2] BERGLUND, J.F., JUNGHENN, H.D. and MILES, P.: Compact
right topological semigroups and generalizations of almost
periodicity, Springer, 1978.
AMS 1980 Subject Classification: 43A60, 42A71
BOHR - FAVARD INEQUALITY - An inequality
appearing in a problem of H. Bohr [1] on the bounded-
ness over the entire real axis of the integral of an
almost-periodic function. The ultimate form of this ine-
quality was given by J. Favard [2]; the latter materially
supplemented the studies of Bohr, and studied the arbi-
trary periodic function
00
I(x) = ~ (ak coskx +bk sinkx)
k =n
with continuous derivative fr)(x) for given coastants r
and n which are natural numbers. The accepted form of
the Bohr- Favard inequality is
lillie ~ K Ilf') lie,
II I II e = XE[O.2"J
max II (x) I,
with the best constant K=K(n, r):
K = Ilf~Pc"1 II I Ik
The Bohr- Favard inequality is closely connected with
the inequality for the best approximations of a function
and its roth derivative by trigonometric polynomials of
an order at most n and with the notion of
Kolmogorov's width in the class of differentiable func-
tions (d. Width).
References
[1] BOHR, H.: 'Un theoreme general sur l'integration d'un
polynome trigonometrique', c.R. Acad. Sci. Paris Ser. I Math.
200 (1935), 1276-1277.
[2] FAVARD, 1.: 'Sur les meilleurs procedes d'approximation de
certaines classes des fonctions par des polynomes
trigonometriques', Bull. Sci. Math. (2) 61 (1937), 209-224; 243-
256.
[3) AmEZER, N.!. [N.J. AKmEZER]: Vorlesungen fiber Approxima
tionstheorie, Akademie-Verlag, 1953 (translated from the Rus-
sian).
LV. Talkov
AMS 1980 Subject Classification: 26D05
BOKS INTEGRAL - One of the generalizations of
the Lebesgue integral, proposed by A. Denjoy (1919)
and studied in detail by T.J. Boks (1921). A real-valued
function f on a segment [a, b] is periodically extended
(with period length b - a) to the entire straight line. For
BOLTZMANN DISTRIBUTION
an arbitrary subdivision a=xo<xl<'" <xn=b of
[a, b], for an arbitrary selection of points I = {~;} 7,
~i E[Xi -I, x;], and an arbitrary t, the following sum is
constructed:
n
I(t) = ~/(~i+t)[Xi-Xi-d.
i=l
If, for p = maxi (Xi - Xi -I )~O, I (t) converges in measure
to a definite limit I, the number I is said to be the Boks
integral (B-integral) of f over [a, b]. Thus, the Boks
integral is an integral of Riemarm type and is a gen-
eralization of the Riemann integral.
The Boks integral represents a considerable extension
of the Lebesgue integral: Any summable function is B-
integrable and these integrals coincide, but there exist
non-summable B-integrable functions; in particular, if g
is the function conjugate with a summable function f,
then it is B-integrable and the coefficients of the series
conjugate with the Fourier series of f are the coeffi-
cients of the respective Fourier series (in the sense of
B-integration) of g (A.N. Kolmogorov). The theory of
the Boks integral was not further developed, since the
A-integral proved to be more convenient for the
integration of functions conjugate with summable func-
tions.
References
[1] BOKs, T.J.: 'Sur les rapports entre les methodes de l'integration
de Riemann et de Lebesque', Rend. Cire. Mat. Palermo (2) 45
(1921),211-264.
[2) ZYGMUND, A: Trigonometric series, 1-2, Cambridge Univ.
Press, 1979.
I.A. Vinogradova
AMS 1980 Subject Classification: 26A39
BOLTZMANN DISTRIBUTION - The statistical
equilibrium distribution function f(p, r) of the momenta
p and coordinates r for the particles of an ideal gas, the
,,+ .fm: l
molecules of which obey the laws of classical mechan-
ics, in an external potential field:
ftp,,) ~ A
U
T (') (1)
Here k is the Boltzmann constant (a universal constant:
k = 1.38 times 1O- 16 erg / degree), T is the absolute
temperature, .; /2m is the kinetic energy of the parti-
cle, U(r) is the potential energy of the particle in the
field, and the constant A is defined by normalization
over a dimensionless phase volume:
j jl(P, r) d 3: :3r = N.
Here N is the total number of particles, h is the Planck
constant (a universal constant
h = 6.62 X 10- 27 erg X sec),
d 3 p = dpx dPy dpzo d 3r = dx dy dz.
415
BOLTZMANN DISTRIBUTION

A can also be defined by the condition of normaliza-
tion in the space of velocities and coordinates, which is
more usual in the kinetic theory of gases:
j jf(P, r)d3vd3r = N,
where t:.U= U(rl)- U(r2). A particu1ar case of (4)
yields the barometric formula, which defines the particle
densities in the gravity field above the surface of the
Earth:
(5)

v =.R..., d 3 v = dvx dvy dvz where g is the acceleration of gravity, m is the mass of
m
the particle, z is the altitude above the Earth's surface,
The Boltzmann distribution is a consequence of the and no is the density at z =0.
Boltzmann statistics for an ideal gas, and is a particular The Boltzmann distribution of a mixture of several
case of the Gibbs distribution gases with different masses shows that the partial den-
nUl!. rio .Pt;. rN) sity distributions of the particles for each individual
P(PI, rl, ... ,PN, rN) = ze
1 kT
component is independent of that of other components.
For a gas in a rotating vessel, U(r) is the field of the
for an ideal gas, when
centrifugal forces:
PT + ~U(ri)'
= ~, 2m
H
, U(r) = _ m~2r ,
and the canonical Gibbs distribution becomes the pro- where w is the angu1ar velocity of rotation.
duct of the Boltzmann distributions for individual par- For references, see Boltzmann statistics.
ticles. The Boltzmann distribution is the limiting case D.N. Zubarev
of quantum statistics for an ideal gas at sufficiently
AMS 1980 Subject Classification: 82A40
high temperatures, when quantum effects can be
neglected. The average occupation number of the i-th BOLTZMANN EQUATION - An equation in kinetic
quantum state of a particle is gas theory, proposed by L. Boltzmann for the determi-
(2) nation of the single-particle distribution function of an
ideal mono-atomic gas [1]. In dimensionless variables
where f.j is the energy corresponding to the i-th quan-
the equation has the form:
tum state of the particle and p. is the chemical potential
defined by the condition ~;;=N F'lrmu1a (2) is valid Ma +(v, 'Yxj)+(F, 'Y.j)
t
= l.L(f,j).
(
(*)
for temperatures and densities at which the average dis-
tance between the particles is larger than the ratio Here f (x, v, t) is the density of the distribution func-
between the Planck constant h and the modulus of the tion of the number of particles in the phase space x v,
average thermal velocity x is the three-dimensional space coordinate, v is the
velocity, t is the time, F is the field strength of the
[~ f/3 lmu. external forces, and f. is a dimensionless parameter
(which is proportional to the ratio between the average
distance covered by the particles between successive
The MaxweU distribution is a special case of the collisions and the typical scale of the phenomena under
Boltzmann distribution (1) for U=O: consideration). In the simplest case the collision opera-
3/2
f() = N [ ~) -p'/2mkT
(3)
tor has the form:
P V 2wkT e .
L(f,j) = j[f(v')j(v'I)-f(v)f(vI)] I v -VI I dwdv lo
The distribution function (I) is sometimes referred to
where v I and v are the velocities of the molecules
as the Maxwell-Boltzmann distribution, the term
before collision, V'I and v' are the velocities of the
Boltzmann distribution being reserved for the distribu-
molecules after collision, and dw is the solid angle ele-
tion function (1) integrated over all momenta of parti-
ment in the direction of the vector VI -v.
cles representing the density of the number of particles
In deriving the Boltzmann equation it is assumed
at the point r:
(4) that the evolution of the function f (x, v, t) is deter-
where no is the density of the number of particles mined by its value at a given moment of time t and by
corresponding to the point at which U=O. The relative the pairwise collisions between the gas molecules, and
densities of the number of particles at different points that the time of interaction between two gas molecules
depend on the differences between the potential ener- during collision is much shorter than the time during
gies at these points: which they move independently of each other. From
the mathematical point of view the derivation of the
~ = e-I!.U/kT,
n2 Boltzmann equation is based on a certain algorithm

416
 BOLTZMANN H-THEOREM

which constructs the operator L in accordance with the operator Lo is non-positive, <cp, Locp> ";;0, and can be
well-known laws of motion of two gas molecules which written as
Lo(cj = -P(c)<p+Gcj>,
collide with one another.
In equation (*) the range of variation of the variable where v(c) (which is sometimes called the collision fre-
t is the half-line t ~o; the range of variation of v is the quency) is a multiplication operator acting on cp and G
entire space R3; and the range of variation of x is a is a completely-continuous integral operator. The func-
subspace n in R3 (n may coincide with R 3). In accor- tion v( c) and the kernel of G have the following form
dance with its physical meaning, the function f (x, v, t) for the hard-sphere model [1]:
should be non-negative and such that
j!(x, v, t)v 2 dv < 00.
P(c) = <W [l.e-
2
c2
/ 2 + [I e I+_1_Jl
21 Ie
jCl e - da]
0
a2

'

The simplest boundary condition on an has the form

!(v -2n(n, v), x, t) = !(v, x, t), XEan, vER 3,
G(e, e') = 471" exp { -I e-e' 12 _ Ie 12_1 e' 12 }_
le-e'l 4 4Ie-e'12
where n is the normal to an. There exist a number of
rigorous statements of the Cauchy problem for equa- -271"1 e-e' I exp { _ e2~e'2}
tion (*).
References The existence theorem for the solution of the Cauchy
[I] BOLTZMANN, L.: Lectures on gas theory, Univ. of California problem as t-'?oo and E:-'?O was proved for the linear-
Press, Berkeley, 1964. ized Boltzmann equation, and the dispersion equation
[2] BOGOLYUBOV, N.N.: 'Selected works', Vol. 2, Kiev, 1970
(in Russian). was studied. The equation is principally employed in
[3] CHAPMAN. S. and COWLING, T.G.: The mathematical theory of the molecular acoustics of ideal gases. The equation
nonuniform gases, Cambridge Univ. Press, 1939. yields correct values for the transfer coefficients (viscos-
A.A. Arsen'ev ity, thermal conduction, velocity of sound) and the
Editorial comments. Stokes - Kirchhoff law of ultra-sound absorption.
References References
[A1] CERCIGNANI, c.: Theory and application of the Boltzmann
[I] CARLEMAN, T.: Probleme mathematique de la theorie kimitique
equation, Scottish Acad. Press, 1975.
des gases, 1960 (in Russian; translated from the French).
[A2] CERCIGNANI, C. (ED.): Kinetic theories and the Boltzmann
[2] ARsEN'EV, A.A.: 'The Cauchy problem for the linearized
equation, Springer, 1984.
Boltzmann equation', USSR Comput. Math. Math. Phys. 5, no.
AMS 1980 Subject Classification: 82A40 5 (1965), 116-136. (Zh. Vychisl. Mat. i Mat. Fiz. 5, no. 5
(1965), 864-882)
A.A. Arsen'ev
BOLTZMANN EQUATION, LINEARIZED, in kinetic
gas theory - A linear integro-differential equation which Editorial comments.
approximately describes the evolution of the one- References
particle distribution function of a sufficiently-rarefied [A 1] CERCIGNANI, C.: Theory and application of the Boltzmann
equation, Scottish Acad. Press, 1975.
gas without internal degrees of freedom for small devia- [A2] Kinetic theories and the Boltzmann equation, Springer, 1984.
tions from equilibrium.
This equation, AMS 1980 Subject Classification: 82A40
acj> _ 1
-at
+ <e, 'VA> - -Lo(cj,
( BOLTZMANN H-THEOREM - One of the conse-
is obtained from the Boltzmann equation quences of the kinetic Boltzmann equation, according to
which the function
H. +<e, 'V x! > -_ -Lif,!)
by substituting
at
1
( H(t) = jh(t, r)dr = j!ln!dpdr,
is a non-increasing function of the time t. Here
! = 71"-3/2 e -c 2
+p.e_c 2 /2cj>
f= f(t, r, p) is the dimensionless classical single-particle
and equating the terms in which the parameter J.L distribution function of the coordinates r and the
appears in the first degree. The operator Lo is said to momenta p, which satisfies the Boltzmann equation.
be the linearized collision operator. The linearized The evolution of the density h (t, r) with time is deter-
Boltzmann equation gives a satisfactory description of mined by the relaxational nature of the evolution of the
the evolution of the distribution function only if function f towards the local MaxweU distribution, while
sup I !(x, e, t)_71"-3/2 e -c 2 I 1.
the limit value of the H-function as t-'?oo is equal to
X,c,t the entropy of the ideal gas, calculated according to
If certain very general assumptions are made, the Gibbs, with the opposite sign. If the increments of t

417
BOLTZMANN H-THEOREM

under consideration are much larger than the time probable. The number of different modes of distribu-
needed for the local Maxwell distribution to be esta- tion of N particles over M cells of size Gj containing N j
blished, the quantity (- h (t, r)) may be identified with particles each is
the entropy density, while (- H(t)) may be identified GNi

with the non-equilibrium entropy of an ideal gas. WB ( N ... ) = N! II -'-" N = ~N;.

1.. ; .. M N;. ;
From the point of view of statistical mechanics, the
It is assumed that the particles are totally independent,
principal significance of the Boltzmann H-theorem con-
that they are distinguishable and that the state remains
sists in the mathematical expression of the fundamental
unchanged by rearrangements of particles within the
assumptions of macroscopic thermodynamics, accord-
same cell. In Boltzmann statistics this magnitude
ing to which, for example, an isolated system spontane-
defines the statistical weight, or the thermodynamic
ously tends to the state of thermodynamic equilibrium,
probability, of a state (unlike in ordinary probability, it
the process being accompanied by an increase in
is not normalized to one). In computing the statistical
entropy.
weight it is assumed that a rearrangement of identical
H-theorems are statements analogous to the original
particles does not result in a change of state, and for
Boltzmann H-theorem, but formulated for statistical
this reason the phase volume WB must be divided by a
systems of a different or a more general type, including
factor N!: WB
the case of non-ideal and quantified systems.
W("'N"') =N!
-.
The theory was presented by L. Boltzmann in 1872.
References Phases with volume reduction as above are said to be
[I) SOMMERFELD, A.: Thermodynamics and statistical mechanics, generic phases (as distinct from the original specific
Academic Press, 1956 (translated from the German).
phases).
[2) UHLENBECK, G.E. and FORD, G.V.: Lectures in statistical
mechanics, Amer. Math. Soc., 1963. All microscopic states with different particle distribu-
I.A. Kvasnikov tions over phase cells when the number of particles
AMS 1980 Subject Classification: 82A05 M
N = LN;
i=1
BOLTZMANN STATISTICS - Statistics used in a
system of non-interacting particles obeying the laws of and the total energy
classical mechanics (a classical ideal gas). The distribu-
tion of particles of an ideal gas (since they do not ;=1
interact) is not considered in the phase space of all par-
(fj is the energy of the particles in the i-th cell) are
ticles (the r-space) as in Gibbs statistical mechanics (d.
given, correspond to the same macroscopic state.
Gibbs distribution), but in the phase space of coordi-
It is assumed that the distribution of the particles in
nates and momenta of one particle (the JL-space). This
a state of statistical equilibrium is the most-probable
is because for an ideal gas the phase volume is
distribution, i.e. corresponds to the maximum
preserved in the JL-space (a special case of one of the W( ... N j ) for a given number of particles N and
Liouville theorems). energy E. The problem of an arbitrary extremum
According to the Boltzmann statistics, this phase W( ... N j ) for given Nand E yields the following
space is subdivided into a large number of small cells Boltzmann distribution for the average number of parti-
with a phase volume such that each cell contains an
cles in a cell:
even larger number of particles N j , and all possible dis- n = N; = e<p-<,)/kT
tributions of the particles over these cells are con- I Gi '

sidered. The phase volume Gj of the i-th cell is its where k is the Boltzmann constant (a universal con-
volume in the JL-space in h 3 -units, where h is the Planck stant k = 1.38 X 10- 16 erg / degree), T is the absolute
constant (a universal constant temperature, and JL is the chemical potential, defined by
h=6.62X 1O- 27 ergXsec). The meaning of such a
the condition N = ~.Nj. In the particular case of a
dimensionless Gj is the maximum number of possible I

potential field U (r):

micro-states in cell i, since the smallest value of the
product of each pair of coordinates and momenta is pl
(, 2m + U(r,),
equal to h according to quantum mechanics, and the
particle has three degrees of freedom. Boltzmann statistics is a special case of Gibbs statistics
Statistical mechanics is based on the assumption that - the canonical ensemble for a gas consisting of non-
all microscopic states corresponding to a given total interacting particles. The Boltzmann statistics is the
energy and a given number of particles are equally limit case of Fermi - Dirac statistics and

418

Bose- Einstein statistics at sufficiently-high tempera-
tures, when quantum effects can be neglected.
The Boltzmann statistics was proposed by L.
Boltzmann in 1868 - 1871.
References
[I] MAYER, J.E. and GoEPPERT-MAYER, M.: Statistical mechanics,
Wiley, 1940.
[2] SOMMERFELD, A.: Thermodynamics and statistical mechanics,
Academic Press, 1956 (translated from the German).
[3] ScHRODINGER, E.: Statistical thermodynamics, Cambridge
Uniy. Press, 1948.
[4] FOWLER, R. and GUGGENHEIM, E.: Statistical thermodynamics,
Cambridge Uniy. Press, 1960.
D.N. Zubarev
AMS 1980 Subject Classification: 80A 15
BOLZA PROBLEM - One of the main problems in
the classical calculus of variations on the conditions for
an extremum, given some contraints on the type of
equations; formulated in 1913 by O. BoIza. The prob-
lem consists of the minimization of the functional
I,
J(x) = jJ(/, x, x)dl+g(tJ, X(tl), 12, X(t2)),
I,
J:RXRnXRn ~ R, g:RXRnXRXRn ~ R,
in the presence of differential constraints on the type of
the equation:
.p(/,x,x)=O, q,:RXRnXRn~Rm, m<m,
and the boundary conditions
!J1h, X(/l)' 12, X(/2)) = 0, 1/1: RXRn XRXRn ~ RP,
P ..;;; 2n+2.
lem, if f
If g-O, the problem is known as the Lagrange prob-
and p<2n +2 it is known as the Mayer
problem. A special feature of the BoIza problem is the
mixed nature of the functional, which is represented as
the sum of an integral functional and a function on the
boundary. In principle, the BoIza problem is equivalent
to the Lagrange problem, and in fact leads to it if one
sets
I,
Jl(x) = j(f(/, x, x)+xn+l)dl,
I,
Xn+l = 0, Xn+l(tl) = ( g ) implies that at least one segment obtained by halving
12 -II
It also leads to the Mayer problem if one sets
J 2 = XO(t2)+g(tJ, x(tl), 12, X(t2)),
where
Xo = J(/, x, x), XO(tl) = 0.
The choice of a given form of the problem, and the
choice of the topology within which the problem is then
considered, is determined by convenience or by expec-
tion in each particular case. In the theory of optimal
control the problem is usually considered in the form
BOLZANO - WEIERSTRASS SELECTION PRINCIPLE
of Mayer's problem; in the classical calculus of varia-
tions it is considered in the form of Lagrange's prob-
lem. The most extensively used topology is that of the
space C 1 of continuously-differentiable functions. To
obtain necessary or sufficient conditions for the
extremum, smoothness requirements must be imposed
on the functions f and g involved in the definition of
the problem and on the mappings q, and l/I, as well as
requirements concerning the regularity of these map-
pings, i.e. the matrices
(al/l / at 1> al/l / ax 1, al/l / at 2, al/l / ax 2) and (aq, / ax)
must have maximal rank p and m, respectively. For the
vector function x(t) to give an extremum in the BoIza
problem (and, similarly, in the Lagrange or Mayer
problem), it is necessary that is satisfies also the Euler
equation and the Weierstrass conditions (for a varia-
tional exfremmn) with respect to the Lagrange function
constructed using the data of the problem with
Lagrange multipliers, as well as the Jacobi condition
and the transversaiity condition.
The notation employed in the above formulation of
the BoIza problem is that used in optimal control
theory. In the classical calculus of variations the BoIza
problem is formulated using a different notation:
x,
J(y) = jJ(x,y,y')dx+g(XJ,Y(Xl),X2,Y(X2)),
XI
References
[I] BLISS, G.A.: Lectures on the calculus of variations, Chicago
Uniy. Press, 1947.
I.B. Vapnyarskir
AMS 1980 Subject Classification: 49BXX
BOLZANO- WEIERSTRASS SELECfION PRIN-
CIPLE - A method of proof which is frequently
employed in mathematical analysis and which is based
on successive subdivision of a segment into halves, after
which the segment having some property is chosen as
the new, initial segment. This method may be
employed if the nature of the segments is such that the
fact that the property is present in some segment
the original segment will also have this property. For
instance, if the segment contains infinitely-many points
of some set, or if some function is not bounded on the
segment, or if a non-zero function assumes values of
opposite sign at the two ends of the segment - all
these are properties of this type. The
BoIzano - Weierstrass selection principle can be used to
prove the Bolzano - Weierstrass theorem and a number
of other theorems in analysis.
Depending on the criterion according to which the
419
BOLZANO - WEIERSTRASS SELECTION PRINCIPLE
segments are chosen in applying the
Bolzano - Weierstrass selection principle, the process
obtained is effective or ineffective. An example of the
former case is the application of the principle to prove
that for a continuous real function that assumes values
of opposite sign at the ends of a given segment, this
segment contains a point at which the function van-
ishes (cf. Cauchy theorem on intermediate values of
continuous functions). In this case the criterion chosen
for the successive choice of the segments is that the
function assumes values of different sign at the two
ends of the chosen segment. If there is a way of com-
puting the value of the function at every point, then,
after performing a sufficient number of steps, it is pos-
sible to obtain the coordinates of the point at which the
function vanishes, to within a given degree of accuracy.
Thus, in addition to proving that a root of the equation
f(x)=O exists on a segment at the ends of which the
values of the function are of different sign, one also has
a method of approximately solving this equation. An
example of an ineffective process is the use of the
Bolzano - Weierstrass selection principle to prove that
a continuous real function on a segment attains a max-
imum on the segment. Here, the segment chosen in the
successive subdivisions is the one on which the max-
imum of the values of the function is not less than that
on the other one. If, as in the former case, it is possible
to calculate the value of the function at any point, this
is still not sufficient for an effective choice of the
required segment. Accordingly, the
Bolzano - Weierstrass selection principle can be used in
this case only to prove an existence theorem which says
that the function assumes its maximum at some point,
but not to specify this point within a given degree of
accuracy.
There exist various generalizations of the
Bolzano- Weierstrass selection principle, e.g. to apply
it in the n-dimensional Euclidean space (n = 2, 3, ... )
to n-dimensional cubes, which are successively subdi-
vided into congruent cubes with side-lengths of one-
half that of the original cube.
L.D. Kudryavtsev
AMS 1980 Subject Classification: 26A03
BOLZANO - WEIERSTRASS THEOREM - Each
bounded sequence of numbers contains a convergent
subsequence. The theorem applies both to real and
complex numbers. It can be generalized to include
more-general objects, e.g. any bounded infinite set in
n-dimensional Euclidean space has at least one limit
point in that space. There exist analogues of this
theorem for even more-general spaces.
The theorem was demonstrated by B. Bolzano [1]; it
420
was later also independently deduced by K. Weier-
strass.
References
[1] BOLZANO, B.: Abh. Biochem Ges. Wiss. (1917).
L.D. Kudryavtsev
AMS 1980 Subject Classification: 26A03
BONNESEN INEQUALITY - One of the more pre-
cise forms of the isoperimetric inequality for convex
domains in the plane. Let K be a convex domain in the
plane, let r be the radius of the largest circle which can
be inserted in K, let R be the radius of the smallest cir-
cle containing K, let L be the perimeter and let F be
the area of K. The Bonnesen inequality [1]
Il = L2_4'lTF;;;' .,r.(R -r)2
is then valid. The equality A=O is attained only if
R = r, i.e. if K is a disc. For generalizations of the Bon-
nesen inequality see [2].
References
[I] BONNESEN, T.: 'Ueber eine Verschlirferung der isoperimetische
Ungieichheit des Kreises in der Ebene und auf die
Kuge10berflache nebst einer Anwendung auf eine Min-
kowskische Ungieichheit fUr konvexe K6rper', Math. Ann. 84
(1921), 216-227.
[2] DISKANT, V.I.: 'A genera1ization of Bonnesen's inequalities',
Soviet Math. Dokl. 14, no. 6 (1973), 1728-1731. (Dokl. Akad.
Nauk SSSR 213, no. 3 (1973), 519-521)
A.B. Ivanov
AMS 1980 Subject Classification: 52A 10, 52A40
BONNET NET - An isothermal net whose lines have
constant geodesic curvature. The square of the line ele-
ment in the parameters of this net is:
du 2 +dv 2
ds 2 = -==----'~
(u+V)2 '
where U= U(u), V= V(v). Studied by O. Bonnet in
1848.
A.B. Ivanov
AMS 1980 Subject Classification: 53A05, 53820
BONNET THEOREM - 1) Bonnet's theorem on the
existence and the uniqueness of a surface with given first
and second fundamental forms [1]. Let the following
two quadratic forms be given:
E du 2 +2F du dv +G dv 2 ,
L du 2 + 2M du dv + N dv 2 ,
the first one of which is positive definite, and let the
coefficients of these forms satisfy the Gauss equations
(cf. Gauss theorem) and the Peterson-Codazzi equa-
tions. Then there exists a surface, which is unique up to
motions in space, for which these forms are, respec-
tively, the first and the second fundamental forms.
2) Bonnet's theorem on the diameter of an oval sur-

face: If the curvature of an oval surface is larger than
or equal to 1 / A 2 at all its points, then the external
diameter of this surface is smaller than 7TA; this esti-
mate cannot be improved. Stated by O. Bonnet in 1855.
A.B. Ivanov
Editorial comments. A proof of this theorem of Bonnet
may be found in [A1] or [A2]. The Peterson-Codazzi equa-
tions are usually called the Mainardi- Codazzi equations, cf
[A1], after G. Mainardi (1857) and D. Codazzi (1868), who
established them.
References
[A1] BLASCHKE, W. and LEiCHTWErB, K.: Elementare Differential-
geometrie, Springer, 1973.
[A2] CARMO, M. DO: Differential geometry of curves and surfaces,
Prentice-Hall,1976.
3) Bonnet's theorem on the mean value, second mean-
value theorem [2r. Let f (x), Ij>(x) be integrable func-
tions on a segment [a, b] and let Ij>(x) be a positive
decreasing function of x; then there exists a number
in [a, b] for which the equality
b ~
jf(x)<p(x)dx = <f>(a)jf(x)dx
a etc.
is valid. If Ij>(x) is merely required to be monotone,
Bonnet's theorem states that there exists a point ~ in
[a, b] such that
b ~ b important:
jf(x)<p(x)dx = <f>(a)jf(x)dx+<f>(b)jf(x)dx
a a ~
is true.
References
[lA] BONNET, 0.: J. Ecole Poly technique 24 (1865), 204-230.
[IB] BONNET, 0.: J. Ecole Poly technique 25 (1867), I-lSI.
[2] BONNET, 0.: 'Remarques sur quelques integrales definies', J.
Math. Pures Appl. 14 (1849), 249-256.
T. Yu. Popova
Editorial comments. Bonnet's original article is [A1].
References
[Al] BONNET, 0.: C.R. Acad. Sci. Paris 40 (1855),1311-1313.
AMS 1980 Subject Classification: 53A05, 26A24,
53C45
BOOLEAN ALGEBRA, Boolean lattice - A partially
ordered set of a special type. It is a distributive lattice
with a largest element '1', the unit of the Boolean alge-
bra, and a smallest element '0', the zero of the Boolean
algebra, that contains together with each element x also
its complement - the element Cx, which satisfies the
relations
sup{x, Cx} = 1, inf{x, Cx} = O.
The operations sup and inf are usually denoted by the
symbols V and 1\, and sometimes by U and
respectively, in order to stress their similarity to the
set-theoretical operations of union and intersection. The
BOOLEAN ALGEBRA
notation X, x' or -x may be employed instead of Cx.
The complement of an element in a Boolean algebra is
unique.
A Boolean algebra can also be defined in a different
manner. Viz. as a non-empty set with the operations C,
V, 1\ which satisfy the following axioms:
1) xVY = Y V x, x 1\Y = Y I\x ;
2) xV(yVz)=(xVy)Vz, x I\(yl\z)=(xl\y)l\z ;
3) (xl\y)Vy=y, (xVy)l\y=y;
4) xl\(yVz)=(xl\y)V(xl\z),
x V(y I\z)=(x Vy)l\(x V z);
5) (xI\Cx)y=y, (xVCx)l\y=y.
If this approach is adopted, the order is not assumed
to be given in advance, and is introduced by the fol-
lowing condition: x~ if and only if x=xl\y.
Other axiomatics are also possible. The axioms of a
Boolean algebra reflect the analogy between the con-
~ cepts of a 'set', an 'event' and a 'statement'. The
'order' relation in a Boolean algebra can be interpreted
in various ways - as set-theoretical inclusion, as causal
follow-up of events, as logical follow-up of statements,
In addition to the basic operations C, V, 1\, other
operations in a Boolean algebra can be defined; among
these the symmetric difference operation is particularly
x +2.Y = (x I\Cy)V(v I\Cx).
Alternative notations are xLly, I x -y I
Any Boolean algebra is a Boolean ring with a unit
element with respect to the operations of 'addition'
(+2) and 'multiplication' (1\); any Boolean ring with a
unit element can be considered as a Boolean algebra.
Boolean algebras first arose in the studies of G.
Boole [1], [2] as a tool of symbolic logic. They subse-
quently found extensive application in other branches
of mathematics - in probability theory, topology,
functional analysis, etc. The applications of Boolean
algebras to logic are based on the interpretation of the
elements of a Boolean algebra as statements (d. Alge-
bra of logic), the complement Cx being interpreted as
the negation of the statement x, and the operations 1\
and V as conjunction and disjunction, respectively.
Closely related to logic is another field of application of
Boolean algebras - the theory of contact schemes (cf.
Contact scheme). Boolean algebras are used in the
foundations of probability theory. A field of events, as
studied in probability theory, is a Boolean algebra; here
the inequality x ~y means that an event y follows from
an event x; the '1', the '0' and the Boolean operations
n V, 1\ and C are interpreted corespondingly.
An example of a Boolean algebra is the system of all
subsets of some given set Q partially ordered by inclu-
421
BOOLEAN ALGEBRA
sion. Such a Boolean algebra is denoted by 2Q; its zero
is the empty set, and its unit is the set Q itself. The set
Q \ x is the complement of an element x; the Boolean
operations V and 1\ coincide with union and intersec-
tion, respectively.
Instead of the subsets of Q it is convenient to con-
sider their characteristic functions. The system XQ of
all such functions, with the natural order, is a Boolean
algebra, which is isomorphic to the Boolean algebra 2Q.
In such a Boolean algebra the operations V, 1\, C, and
+2 are interpreted as follows:
(xVyXq) = max{x(q),y(q)},
(x;\y)(q) = min{x(q),y(q)} = x(q)y(q),
(Cx)(q) = 1-x(q),
(x+1Y)(q) = I x(q)-y(q) I = x(q)+y(q) (mod 2), (qEQ).
The following cases are especially important:
1) Q=Qn={1, ... ,n}.
In this case the characteristic functions of the subsets
are 'two-valued symbols' of the form:
O,
x = (XI> ... ,xn ), Xj= { 1,
which are 2n in number. For n = 1, the two-element
Boolean algebra, consisting only of '1' and '0', is
obtained.
2) Q=XQ
In this case, all possible functions, defined on the
system of all binary symbols of length n, and taking the
values '0' and '1' only, are elements of X Q They are
called Boolean functions in n variables (cf. Boolean
function). Every well-constructed formula of predicate
logic defines some Boolean function; if two functions
are identical, the formulas are equivalent.
Under certain conditions a subset E of elements of a
Boolean algebra X is itself a Boolean algebra with
respect to the order induced from X. This is the case,
in particular, if:
(a) E is a principal ideal, i.e. a set of the form
{XEX: x:s;;;u}; the element u then acts as the unit '1';
(b) E is a subalgebra of a Boolean algebra X. This
means that if x, y EE, it follows that
x Vy, X /\y, Cx EE. The '0' and '1' of the initial
Boolean algebra are also the '0' and '1' of the new
Boolean algebra. The Boolean subalgebras of 2Q are
especially important; they are called algebras of sets.
Any set E eX generates a certain subalgebra - the
smallest sub algebra that contains E.
Homomorphisms of Boolean algebras play a special
role under the mappings of Boolean algebras; they are
mappings which commute with the Boolean operations.
A bijective homomorphism of Boolean algebras is an
isomorphism. If a Boolean algebra X is generated by a
set E, then all mappings of E into an arbitrary Boolean
422
algebra have an extension to a homomorphism if and
only if E is an independent set, i.e. if all elements of the
form
XI;\ ... ;\xp;\CXp+l;\ ... ;\Cxm , xjEE, Xj=l=Xk,
are non-zero. A Boolean algebra generated by an
independent system is called a free Boolean algebra.
An example of a free Boolean algebra is the algebra
of Boolean functions in n variables considered above.
Independent generators of it are the functions
J;:J;(x) = J;(XI> ... ,xn ) = Xj.
Stone's theorem: Every Boolean algebra X is iso-
morphic to some algebra of sets, namely, the algebra of
all open-and-c1osed sets of a totally-disconnected com-
pactum l(X), defined up to a homeomorphism. This
compactum is known as Stone's compactum. To a
homomorphism of a Boolean algebra X into a Boolean
algebra Y corresponds a topological imbedding of l( y)
into l(X); to a subalgebra of a Boolean algebra X
there corresponds a continuous image of l(X). The
Stone compactum of a free Boolean algebra is a dyadic
discontinuum.
A Boolean algebra X is called complete if any set
E e X has an upper bound sup E and a lower bound
inf E. This is equivalent to l(X) being extremal (cf.
Extremally-disconnected space). Sub algebras of a com-
plete Boolean algebra containing the bounds of all their
subsets calculated in X are called regular subalgebras.
The weight of a Boolean algebra X is the lowest cardi-
nality of a complete generating set, i.e. a set which is not
contained in any regular subalgebra other than X. If
the weights of all non-zero principal ideals are identi-
cal, then the Boolean algebra is called uniform; such
algebras invariably contain a complete generating
independent set. In other words, a complete uniform
Boolean algebra can be 'stretched onto' a free Boolean
algebra. The study of an arbitrary Boolean algebra
readily reduces to the study of uniform Boolean alge-
bras. An incomplete Boolean algebra can be completed
in different ways, i.e. it can be imbedded as a subalge-
bra in some complete Boolean algebra.
A complete Boolean algebra is called normed if a
real-valued function J.L (a measure) is defined on it with
the following properties: 1) if x=O, then J.L(xO; 2) if
EeX and xl\y=O when x,YEE, x=y, then
p.(supE) = ~p.(x).
XEE
In probability theory, in which normed Boolean alge-
bras are particularly important, it is usually assumed
that J.L( I) = 1. Here, the value of J.L{x) is interpreted as
the probability of an event x. The classical theory of
measure and integral can largely be applied to normed
Boolean algebras. Normed Boolean algebras have been
completely classified [4], [5], [7]. In particular, for uni-

form normed Boolean algebras the only invariant is the
weight. Not all Boolean algebras can be normed. Many
conditions for the existence of a measure are known,
but these are far from exhaustive in the problem of
nonning.
A Boolean algebra can be endowed with various
topologies. A specially important one is the so-called
(0 )- topology, which, for a normed Boolean algebra, IS
metrizable, and corresponds to the metric
p(x,y) = JL(x/\Cy)V(Cx/\y)],
and is identical with the Tikhonov topology for
Boolean algebras of the form 2Q. In the most general
case there need not be a topology compatible with the
order in a Boolean algebra.
References
[I] BOOLE, G.: The mathematical analysis of logic: being an essay
towards a calculus of deductive reasoning, Macmillan, 1847.
[2] BOOLE, G.: An investigation of the laws of thought, on which are
founded the mathematical theories of logic and probabilities,
Dover, reprint, 1951.
[3] SIKORSKI, R.: Boolean algebras, Springer, 1969.
[4] VLADIMIROV, D.A.: Boolesche Algebren, Akademie Verlag, 1978
(translated from the Russian).
[5] liALMOS, P.R.: Lectures on Boolean algebras, v. Nostrand,
1963.
[6] RASIOWA, E. and SIKORSKI, R.: The mathematics of
metamathematics, Polska Akad. Nauk., 1963.
[7] STONE, M.H.: The theory of representations for Boolean alge-
bras', Trans. A mer. Math. Soc. 40 (1936), 37-111.
[8] BIRKHOFF, G.D.: Lattice theory, Colloq. Pub!., 25, Amer. Math.
Soc., 1973.
[9] HERMES, H.: Einfiihrung in die Verbandstheorie, Springer, 1967.
[10] KOLMOGOROV, A.N.: 'Algebres de Boole metriques completes',
in VI Zjazd Matematykow Polskich, Krakow, 1950.
[II] DUNFORD, N. and SCHWARTZ, J.T.: Linear operators. Spectral
operators, 3, Interscience, 1971.
[12A] KAKUTANI, S.: 'Concrete representations of abstract (L)-spaces
and the mean ergodic theorem', Ann. of Math. (2) 42, no. 2
(1941),523-537.
[12B] KAKUTANI, S.: 'Concrete representations of abstract (M)-
spaces (a characterization of the space of continuous func-
tions)" Ann. of Math. (2) 42, no. 4 (1941),994-1024.
[13] MAIiARAM, D.: Proc. Nat. Acad. Sci. USA 28 (1942), 108-111.
[14] MACKEY, G.W.: The mathematical foundations of quantum
mechanics, Benjamin, 1963.
[15] YOSIDA, K.: Functional analysis, Springer, 1980.
[16] KURATOWSKI, K.: Topology, 2, Acad. Press, 1966-1968
(translated from the French).
D.A. Vladimirov
AMS 1980 Subject Classification: 06EXX
BOOLEAN EQUATION - An equation of the form
!(Xj, ... ,xn ) = 0, (*)
where f is a Boolean function in n variables. The set of
all solutions of an equation of the form (*) can be
described by a system of Boolean functions depending
on n arbitrary parameters.
References
[I] BIRKHOFF, G.: Lattice theory, Colloq. Pub!., 25, Amer. Math.
Soc., 1973.
T.s. Fofanova
AMS 1980 Subject Classification: 06E30
BOOLEAN FUNCTION
BOOLEAN FUNCI10N, function of the algebra of
lOgiC - A function whose arguments, as well as the func-
tion itself, assume values from a two-element set (usu-
n.
ally {O, I Boolean functions are one of the main sub-
jects of discrete mathematics, in particular, of
mathematical logic and mathematical cybernetics.
Boolean functions first occurred in the mathematical
formulation of logical problems, and were named after
G. Boole, who laid the foundation for the applications
of mathematics in logic in the middle of the nineteenth
century; cf. Algebra of logic.
One such problem is the construction of an algebra
of propositions. To this purpose, one of the two values
o ('false') or I ('true') is assigned to each proposition;
the principal logical relations 'and', 'or', 'not', 'if ...
then', etc., can then be regarded as the respective 'ele-
mentary' Boolean functions: xl\y, xVy, ex, x~y, etc.
When this is done, the value of any complex proposi-
tion, constructed with the aid of the fundamental logi-
cal connectives from given propositions, is a Boolean
function of the values of these propositions. Such a
Boolean function is a composition of elementary
Boolean functions corresponding to the logical connec-
tives forming part of the complex proposition. It
became clear later that the language of Boolean func-
tions is suited for a description of the operation of
discrete control systems (cf. Control system) such as
contact schemes, diagrams of functional elements, logi-
cal networks, switching networks, etc. These control
systems are constructed in accordance with certain
rules from a number of initial elements, just as complex
statements are constructed from simple ones. The rules
governing the construction of such control systems as
well as the operation of the initial elements are such
that the operation of complex control systems can be
described with the aid of Boolean functions. Boolean
functions are also used in certain problems of integer
programming that reduce to solving a system of
Boolean equations of the form
!l(Xj, ... ,xn ) = ... = !m(xj, ... ,xn ) = 0,
where the j;, i = I, ... ,m, are Boolean functions.
There are also other ways of using Boolean functions in
discrete mathematics, so that the study of Boolean
functions is of interest in its own right.
An essential feature in solving problems related to
Boolean functions is the method of specifying Boolean
functions. There are several such methods: tables, for-
mulas, special classes of formulas known as normal
forms (cf. Boolean functions, normal forms of), subsets
of the vertices of an n-dimensional unit cube, etc. In
the last-named case every selection of length n of values
of the arguments (0 or I) is regarded as a vertex of the
n-dimensional unit cube, in which case a Boolean func-
423
BOOLEAN FUNCTION
tion of n arguments can be defined by the subset of the
comers at which it assumes the value 'I'. This subset,
when written out as a matrix whose rows are selections
of values of the arguments of the Boolean function, is
known as a Boolean matrix. If a Boolean function
describes the operation of control systems, the latter
can also be regarded as a method of specifying the
Boolean function. One usually says that this control
system realizes the given Boolean function. The realiza-
tion of Boolean functions by many kinds of control
systems is closely related to a large number of prob-
lems, such as synthesis, minimization, control and relia-
bility problems, etc. Another type of problem arises in
the study of properties and classes of Boolean functions
specified by different methods; this is the study of the
metric characterizations of various classes of normal
forms of Boolean functions and the related geometrical
properties of the n-dimensional unit cube (cf. Boolean
functions, metric theory of), as well as of the various
algebras of Boolean functions (ct. Many-valued logic;
Equivalent transfonnations). The system of all classes of
Boolean functions that are closed under composition
was described by E. Post. It forms a countably-infinite
lattice with five maximal (pre-complete) classes.
In certain cases it may be necessary to consider par-
tial (i.e. not everywhere-defined) Boolean functions for
which the problems listed have a specific character.
References
[1] NOVIKOY, P.S.: Elements of mathematical logic, Edinburgh,
1964 (translated from the Russian).
[2] YABLONSKII, S.V., GAVRILOV, G.P. and KUORYAVTSEV, V.B.:
Functions of the algebra of logic and Post classes, Moscow, 1964
(in Russian).
v'B. Kudryavtsev
Editorial comments. Cf. [A2] for uses of Boolean func-
tions and equations in switching theory and logical net-
works. Applications of Boolean functions and equations in
operations research are discussed in [A1].
References
[A1] HAMMER, P.L. and RuoEANU, S.: Boolean methods in opera-
tions research, Springer, 1968.
[A2] RuoEANU, S.: Boolean functions and equations, North-
Holland,1974.
AMS 1980 Subject Classification: 06E30
BOOLEAN FUNCfIONS, METRIC THEORY OF -
The branch of mathematics in which one studies
numerical characteristics and metric properties of
Boolean functions. The principal parts of this theory
are concerned with properties of 'almost-all' Boolean
functions (cf. Boolean functions, minimization of), pro-
perties of the set of all Boolean functions in a given
number of variables, and special subclasses of Boolean
functions. Another topic is the structure of the truth
domains of Boolean functions with the aid of numerical
characteristics that appear mainly in problems con-
424
nected with the minimization of Boolean functions and
the theory of local algorithms. These characteristics are
the dimension and the extent of functions.
Let Nf(x" ... ,x,) be the set of vertices of the n-
dimensional unit cube on which a function
f(x], ... ,xn ) is equal to one. Consider all maximal
intervals of the function f(x], ... ,xn ) and select the
interval of the largest dimension r. The quantity r is
called the dimension of f (x], ... ,xn ) and is denoted
by Dimf(XI, ... ,xn ). In terms of the dimension one
can estimate the complexity ratio of the most complex
dead-end and the shortest disjunctive normal forms of f
(cf. Boolean functions, normal forms of). An upper esti-
mate of this ratio is 2Dimf. The inequality
Dimf(x\> ... ,xn ) E;;; [lo~nl+l
is valid for 'almost-all' Boolean functions.
In solving problems of minimization of Boolean
functions it is of interest to calculate the dimension of
'typical' maximal intervals. It has been proved that
'almost-all' maximal intervals of 'almost-all' Boolean
functions f(x], .. . ,xn ) are of dimension close to
lo&lo~n.
Let Sk(,iJl,j) be the k-th order smooth neighbour-
hood (cf. Algorithm, local) of an elementary conjuction
~ occurring in the abridged disjunctive normal form SRf
of f, and let k (~,f) be the minimal value of the order
of the neighbourhood in which Sk(~' j) includes all
elementary conjunctions occurring in the abridged dis-
junctive normal form SRf . The quantity
p(j) = max k(~, f)
lle91f
is called the extent of f For 'almost-all' Boolean func-
tions
n
p(j(x\> ... ,xn )) ~ 1 1
o~ o~n
Let
p(n) = j(x,max
... ,x,)
P(j(Xlo'" ,xn )).
It is known that p (n) can be realized on a Boolean
function of a special kind, which is called a chain. A
function ll-{X) , . . . ,xn ) is called a chain if the set N%f
its zeros can be represented as a sequence a), ... ,aq
such that p(aj, aj + 1) = 1, where p is the Hamming dis-
tance (cf. Code); the distance between other pairs a" as
(possibly except for the pair (a), aq is larger than one,
and no interval of dimension 2 is completely contained
in the set of ones of 1/;. The extent of the chain 1/; is
q - 1. Therefore, the problem of calculating p (n)
reduces to the construction of a chain with maximal q
in the n-dimensional unit cube. It has been shown by
direct construction of such chains that
Cj2 n <p(n) < C2'2n,
where CJ, C2 are constants. The constmction of closed

chains (cycles), i.e. of chains in which p(a" aq )= 1, with
maximum cardinality of the set N", is an important
part of the proof of the theorem that the properties of
conjuctions 'of occurring in the minimal or shortest dis-
junctive normal form' are non-computable in the class
of local algorithms.
The following result clarifies the structure of the
truth domains of 'almost-all' Boolean functions. A set
M of vertices of the n-dimensional unit cube is called
connected if for any yoint aEM there exists a point {J
a
in M such !!tat p(a, {J) = 1. Ayoint in M is called iso-
pted if all {J such that p(a, {J)= I satisfy the condition
{Jfl.M. The following theorem holds: For 'almost-all'
Boolean functions f (x), ... ,xn ) the set of ones,
Nf(x" ... ,x.), split~ into the sum of a connected set and
a set of isolated points. The cardinality of the con-
nected set is not less than
2n -)-Vn2"/ 2 -constlo!;2n and the number of iso-
lated points does not exceed constl0!;2 n.
The results obtained in calculating the extent of
'almost-all' Boolean functions are closely connected
with the results obtained on calculating the radii and
diameters of the graphs generated by Boolean func-
tions. The graph G(j) generated by a Boolean function
f has as vertices the points of Nf and as edges the pairs
of points of Nf at H~g distance one from each
other,: The distance rG(a, {J) between two vertices a The estimate of m (j) can be improved if supplemen-
and {J of G (j) is defined as the length of the shortest
chain connecting witha p. (It is assumed that the ver-
a p
tices and belong to the same connected component
a
of G(j).) The deviation of the vertex in G(j) is the
quantity l(a)=maxrG(a, P), where the maximum is
taken over all the vertices of G that belong to the same
connected component as a. The radius of a connected
component K of G is the number R(K)=minaEKI(a).
The quantity R(G)=maxR(K), where the maximum is
taken over all connected components of G, is called the
radius of G. The diameter of G is the number
D(G)=maxrG(a, P), where the maximum is taken over
all pairs of vertices a, p that belong to the same con-
nected component. For 'almost-all' Boolean functions
f(x), ... ,XII) the quantities R(G) and D(G) are such
that n -2:s;;;;,R(G(j:s;;;;,n -I and D(G(j=n-1.
Of all the results of computations of numerical
characteristics of individual classes of Boolean func-
tions only those concerning monotone Boolean func-
~ons will be considered. Let a=(a)'" all),
{J=({J) . ~ . {JII) be binary ordered samples. One says
that a:s;;;;,{J if ai:s;;;;,{Ji, i=I, ... ,no A Boolean function is
called monotone if a~p implies that f(a)~f({J). It is
of interest to determine the exact number #..n) of dis-
tinct monotone Boolean functions in the variables
X), ... ,XII' This number is known for only a few small
BOOLEAN FUNCTIONS, METRIC THEORY OF
values of n. The asymptotic equality
n
is also known.
The problem of optimal decoding of a monotone
Boolean function has a large number of applications in
the solution of discrete extremal problems. Consider a
given algorithm A for the computation of a monotone
Boolean function f (x), ... ,XII) at any point a. If it
can be established in the course of the computations
that f(a)=1 for some point a,
then f({J)=1 for all
p~a. If f(a)=O, the function f is known (and van-
ishes) at all points p~a. Therefore, to effect the decod-
ing, i.e. to completely reproduce the function f, the
algorithm A need only compute its values in a certain
set of points. Let m (j) denote the minimum number
of points sufficient to reproduce f completely. Let
m(n) = max m(j).
j(x" .. . ,x.)
The quantity m(n) satisfies the asymptotic equality:
lJ
n 1
m(n) - 2 [;
tary information about the monotone Boolean function
is available.
A problem related to the problem of decoding mono-
tone Boolean functions is that of calculating numerical
characteristics of the set of essential variables of
Boolean functions that are not defined everywhere, i.e.
functions defined on some subset of the set of vertices
of the n-dimensional unit cube. A collection of vari-
ables Xi" ... ,Xi. is called essential for a Boolean func-
tion F(x" ... ,XII) that is not everywhere defined if
{Xi" . . . ,xd C;; { X ), . . . ,XII} and if there exists a
Boolean function CP(Xi" ... ,Xi.) that is not everywhere
defined such that F(x), ... ,XII)-CP(Xi" ... ,Xi.)
throughout the domain of definition of F. An essential
collection is called a dead-end collection for F if no
proper part of it is essential for F.
Each everywhere-defined Boolean function has a
unique dead-end collection. For Boolean functions that
are not defined everywhere the structure of the collec-
tions of essential variables is distinguished by its large
variety. Let ~ be the system of subsets of the set
{x), ... ,xn } that satisfies the following condition: If
S E~ and if S' is an arbitrary subset of {x), ... ,XII}'
then SUS' E~. Then there exists a Boolean function
F~(x" ... ,XII)' not everywhere defined, such that ~ is
a system of essential collections for F~. Let tF(n) be
425
BOOLEAN FUNCTIONS, METRIC THEORY OF
the number of distinct dead-end collections for
F(Xb ... ,XII) and put
t(n) = F(x"max tp(n).
... ,x,,)
It is known that
The algorithm for the construction of all dead-end
collections for F(x., ... ,XII) has maximal computa-
tional labour 2([11 j2) and this maximum is in fact
attained. (Here the unit of labour is defined as the
labour involved in verifying whether or not a given col-
lection of variables is essential for F.)
The metric theory of Boolean functions also deals
with problems of computing characteristics connected
with the problem of minimization of Boolean functions.
References
[I] GLAGOLEV, V.V.: 'Certain estimates of the disjunctive normal
forms of functions of the algebra of logic', Problemy Kibemet.
19 (1967),75-94 (in Russian).
[2] ZlruRAvuv, Yu.l.: 'On algorithmic extraction of the set of
essential variables of not everywhere defined functions of the
algebra of logic', Problemy Kibemet. 11 (1964), 271-275 (in
Russian).
[3] KOROBKOV, V.K.: 'On monotone functions of the algebra of
logic', Problemy Kibemet. 13 (1965),5-28 (in Russian).
[4] SAPOZHENKO, A.A.: 'Metric properties of almost all functions
of the algebras of logic', Diskret. Anal. 10 (1967),91-119 (in
Russian).
[5] SAPOZHENKO, AA.: 'The order of the neighbourhood of maxi-
mal intervals for almost all logical functions', Soviet Math.
Dokl. 9, no. 3 (1968), 591-594. (DokL Akad. Nauk SSSR 180,
no. I (1968), 32-35)
[6] KLBITMAN, D.l.: 'On Dedekind's problem: the number of
monotone Boolean functions', Proc. Amer. Math. Soc. 21
(1969), 677-682.
Yu.I. Zhuravlev
AMS 1980 Subject Classification: 06E30
BOOLEAN FUNCflONS, MINIMIZATION OF - A
representation of Boolean functions by normal forms
(cf. Boolean functions, normal fonus of) that are most
simple relative to some measure of complexity. The
usual meaning of the complexity of a normal form is the
number of letters in it. A simplest form is then called a
minimal form. A measure of complexity sometimes used
is the number of elementary conjunctions in a disjunc-
tive normal form or the number of factors in a con-
junctive normal form. In this case a simplest form is
called a shortest form. In view of the duality of disjunc-
tive and conjunctive normal forms, it is sufficient to
consider disjunctive normal forms only.
The construction of shortest and minimal disjunctive
normal forms each has its own specific features. The
sets of minimal and shortest disjunctive normal forms
of one and the same function may be connected by
set-theoretical relations: of being contained one in one
426
another, of having an empty intersection or a non-
empty symmetric difference. Let mf be the complexity
of the minimal disjunctive normal form of a function f,
let kf be the minimal complexity of its shortest disjunc-
tive normal form and let /(n) be the largest of the
ratios kf / mf over all functions in n variables. Then the
following asymptotic relation holds:
n
I(n) '" 2'
By the problem of minimization of Boolean functions
one usually understands that of constructing their
minimal disjunctive normal forms. There is a trivial
algorithm for constructing all minimal disjunctive nor-
mal forms of an arbitrary Boolean function
f (x 1, . . . , XII)' which operates as follows: All disjunc-
tive normal forms in the variables Xl"'" XII are
reviewed, and those which realize the function f and
have minimal complexity are selected. In fact, this algo-
rithm is not applicable even for small n, since the
number of operations required increases rapidly. Many
other algorithms have therefore been constructed,
which, however, are not effectively applicable to all
functions.
The initial specification of a function in the problem
of minimization is usually a table, a perfect disjunctive
normal form (cf. Boolean functions, normal forms of) or
an arbitrary disjunctive normal form. The first stage
consists in the transition to the so-called abridged dis-
junctive normal form, which is uniquely determined for
each function. Many methods for the realization of this
transition are available. The most universal method
consists in performing transformations on a disjunctive
normal form of the type
A V AB ~ A (absorption),
xA V CxB ~ xA V CxBV AB.
A typical property of the abridged disjunctive normal
form is that it is possible to obtain any minimal
abridged disjunctive normal form from it by removing
certain elementary conjunctions. The second, most
laborious stage is to extract from the abridged disjunc-
tive normal form all dead-end disjunctive normal
forms, among them all the minimal ones. A geometri-
cal representation of Boolean functions is usually
employed at this stage. Let Ell denote the set of all
vertices of the n-dimensional unit cube. Each Boolean
function f (x J, ,XII) is in one-to-one correspondence
with the subset Nf , Nfr;;,En, of vertices a such that
f (a) = 1. Let ~ be an elementary conjunction of rank r.
The set N m: is then called an interval of rank r
corresponding to the elementary conjunction ~. One
says that a system of intervals N m:" ,N m:m forms a
covering of a set N r;;,En if

N = N., U ... UN.~.
Since the equalities
f= ~lV V~m and NI = N., U ... UN~
are equivalent, the problem of minimization of Boolean
functions is equivalent to the search for coverings with
a minimal sum of the ranks of their intervals. Such
coverings are called minimal. An interval N g{ is called
maximal for a function f if N III r;,Nf and if there is no
interval NI1J such that N'I. CN I1J r;,Nf. The construction
of dead-end disjunctive normal forms of a function f
reduces to a search for coverings of Nf by maximal
intervals such that no proper subset of it is a covering
of Nf . These coverings correspond to dead-end dis-
junctive normal forms and are called irreducible. They
are obtained from ~verings corresponding to abridged
disjunctive normal forms by removal of certain inter-
vals.
The selection of the minimal disjunctive normal
forms from all dead-end disjunctive normal forms is
also a very laborious process, since 'almost-all' Boolean
functions in n arguments have no fewer than 22n dif-
ferent dead-end disjunctive normal forms. The diversity
of the complexities of dead-end disjunctive normal
forms may be very wide, so that if a random rather
than a dead-end minimal disjunctive normal form is
selected, a large error may result.
Estimates of the number of dead-end disjunctive nor-
mal forms and of the scatter of their complexities show
that a more detailed inspection during the elimination
of the elementary conjunctions from a disjunctive nor-
mal form is a natural way of improving the effective-
ness of minimization algorithms. A conjunction should
be eliminated (or, on the contrary, retained until the
end of the process) only if it can be established by
some procedure that it does not occur in any minimal
disjunctive normal form for f (occurs in all minimal
disjunctive normal forms for f). The latter fact is usu-
ally established by an analysis of the conjunctions that
are close to the one in question, i.e. occur in a neigh-
bourhood of it (cf. Algorithm, local). Here one has to
accumulate information about elementary conjunctions
and use it in the subsequent analysis. Procedures of this
kind are known as local Simplification algorithms. Some
of them are quoted below, and their description
involves the concept of the k-th order neighbourhood
Sk(~' 91) of an elementary conjunction ~ in a disjunctive
normal form 91, ~e91.
The zero-order neighbourhood So(~, 91) consists of
the single conjunction ~. If Sk -I (~, 91) is the neigh-
bourhood of order k - I, then the neighbourhood
Sk(~' 91) of order k consists of all conjunctions ~j of
the disjunctive normal form 91 that satisfy one of the
following conditions: I) N Ill} has a non-empty intersec-
BOOLEAN FUNCTIONS, MINIMIZATION OF
tion with at least one interval corresponding to a con-
junction of Sk -I (~, 9C); or 2) N'l.Jr;, U ; =1N 113" where
N I1J" i=l, ... ,r, satisfies I).
Quine's algorithm. Here one considers at each step
the neighbourhood S 1 (~, 91) of one of the conjunctions
in the disjunctive normal form 91. As the algorithm is
executed for each conjunction, one attempts to com-
pute one of the following properties:
p 1 (~, 91) - '~ occurs in all minimal disjunctive nor-
mal forms'; and
p 2 (~, 91) - '~ does not occur in any minimal dis-
junctive normal form'.
The algorithm works as follows. I) The neighbour-
hood SI(~' 91)={~, ~I. . . . '~m} is successively
formed for each conjunction ~ in 91. The inclusion
N III C U ~= 1N Ill, is verified. If this is not the case, then
~ is marked in some way as belonging to all minimal
disjunctive normal forms. One says that ~ is part of the
kernel of 91 (is a kernel conjunction). 2) Let the con-
junctions ~I' .. ,)8q be marked in 91 in the first stage.
The remaining conjunctions in 91 are ordered, and the
inclusion N!B r;, U r= 1N 113, is checked for each such
conjunction)8. Conjunctions satisfying this relation do
not occur in any minimal disjunctive normal form and
are removed from 91.
Regular points algorithm. During each step of this
algorithm one examines the neighbourhood S 2(~' 91) of
a conjunction ~ in the disjunctive normal form 91, and
the conjunctions not occurring in any dead-end disjunc-
tive normal form, and hence not occurring in any
minimal disjunctive normal form either, are eliminated.
The description of the algorithm involves the concept
a
of an a-bundle in 91, which, for this point in N~, is
the set M(a,91) of intervals N'l.J such that ~je91,
a
aeN\1lJ' For a conjunction ~ in 91 a point of N\1l is
called regular with respect to (~, 9C) if there exists a
point a' such - that a'eN~ \NIll and
M(a', 91)r;,M(a, 91). A set M r;,N'I. is called regular
with respect to (~, 91) if all its points are regular with
respect to (~, 91). This algorithm is based on the fact
that a conjunction ~ of the abridged disjunctive normal
form of a function f does not occur in any dead-end
disjunctive normal form of f if and only if N'I. is a reg-
ular set with respect to (~, 91). The algorithm checks
whether in a certain sequence the conjunction intervals
occuring in the disjunctive normal form are regular sets
and eliminates those that are regular. Whether or not
an interval N III is regular with respect to (~, 91) is com-
pletely determined by the neighbourhood S2(~' 91) but
not, in general, by the neighbourhood S 1(~, 91).
The A -algorithm. The concept employed here is that
of a disjunctive normal form that is minimal with
respect to a disjunctive normal form 91, i.e. one that is
427
BOOLEAN FUNCTIONS, MINIMIZATION OF
minima] among all disjunctive normal forms that are
equivalent to 9l and are obtained from 9l by omission
of certain elementary conjunctions. Two properties of
an elementary conjunction in a disjunctive normal form
are considered:
PI (m, 9l) - 'm occurs in all disjunctive normal forms
that are minimal with respect to 9l'; and
P 2 (m,9l) - 'm occurs in no disjunctive normal form
that is minimal with respect to 9l'.
It is assumed that Pi = I if the property Pi holds, and
Pi =0 otherwise. It is also assumed that disjunctive nor-
mal forms are formed from conjunctions with informa-
tive marks ("'10 "'2), where "'iE{O, I, a}. The equality
a
"'i :: means that the property Pi has not been calcu-
lated (i = I, 2), while the equalities "'i = I mean that
=
Pi ="'i (i 1,2). A conjunction m
with informative
marks ("'I, "'2) is denoted by m""""'. The A-algorithm
computes the values of the properties P I and P 2 for
the conjunctions min the disjunctive normal form 9l,
using for this purpose only the conjunctions of the
neighbourhood S2(m,9l) and their informative marks.
For a conjunction mof rank r of 9l a set Mr;;.Nrtf. is
called a set of the first type relative to (m, 9l) if 9l con-
tains conjunctions mlo , mm of ranks r., ... ,rm
such that Mr;;. U j=INrtf.j and r> ~j=lrj.
The difference 91 1 \ 912 of two disjunctive normal
forms 911 and 912 is the disjunctive normal form con-
sisting of the elementary conjunctions occurring in 9l]
but not in 9l2
Before the application of the A-algorithm all con-
junctions of the disjunctive normal form 9l have the
mark (a, a). Mter i steps have been performed and the
conjunctions m;, ... ,m~ have received the mark (1, a),
and the conjunctions ~~, ... ,~; have received the
mark (a, I), the (i + I)-th step consists of the following.
The conjunctions in the disjunctive normal form are
ordered in some way. If the disjunctive normal form
9l\(Vj=]~jVVj=]~;) is empty, the A-algorithm is
terminated. If it is not, the conjunctions of this disjunc-
tive normal form are ordered in some way; the first
conjunction ~ in the sequence and its neighbourhoods
S] and S 2 in the disjunctive normal form
~=9l\(Vj=]~;) are singled out, and the relation
N'/I. = U. Nmo
!8e(S, ('/I., 91) \ 'II.)
is checked. If it is satisfied, the mark (a, a) over ~ is
replaced by (1, a) and the (i + l)-th step of the A-
algorithm is terminated. If not, then a check is made
whether or not N rtf. can be represented in the form
M] UM2' where M] is a regular set relative to (m, 9l)
and M 2 is a set of the first type relative to (~, 9l). If
N'l can be represented in this form, the mark (a, a)
over ~ is replaced by (a, 1) and the (i + 1)-th step of
428
the A -algorithm is terminated; if this is not possible,
the procedure is applied to the second conjunction in
the sequence, etc. If the mark does not change in any
one of these conjunctions, then after all the conjunc-
tions have been examined, the operation of the A-
algorithm is terminated.
If the abridged disjunctive normal form 9lf of a
function f is taken for 9l, then no conjunctions that
obtained in the algorithm the mark (a, I) occur in any
minimal disjunctive normal form of f, they are elim-
inated from 9lf . Conjunctions that obtained the mark
(I, a) occur in all minimal disjunctive normal forms of
f Various special cases of the A -algorithm have also
been considered.
The ring algOrithm places over the conjunctions infor-
mative marks ("'10 "'2), with the same meaning as in the
A -algorithm. At each step of the ring algorithm use is
made of the conjunctions contained in the k-th order
neighbourhood of some conjunction and of their infor-
mative marks. A simplified version of this algorithm is
outlined below. The ring algorithm in its complete form
is the best local algorithm with a special memory rela-
tive to the neighbourhoods Sk(m,9l). The algorithms
described above are special cases of the general ring
algorithm. If
Sk-1(~' 91) = {~, ~\o 0 ,~tl,
Sk(~' 91) = {~, ~\o '~/' ~I+\O '~m}
0
and
I m
N s. , = N'/I. U UN'/I." N s = N'/I. U UN'/I.J'
;=1 j=1
Q(Sk) = Ns \ N s. , ,
then to each subset Nr;;.Q(Sk) a Boolean function
f(x 10 ... ,xn ) that is not everywhere defined is
assigned, such that the set M{ of its ones is Ns. \ N
and the set M6 of its zeros is En \ Ns.; the function f is
not defined on N. The set of such functions is denoted
by Fk(m). Prior to the beginning of the algorithm the
conjunctions of the disjunctive normal form 9l in ques-
tion have the mark (a, a). If, after i steps have been
completed, conjunctions ~'I' ... ' with the markm:
(1, a) and conjunctions ~'{, ... ,~; with the mark
(a, I) are obtained, then the (i + l)-th step is performed
as follows. The conjunctions of the disjunctive normal
forms are ordered in some way. If the disjunctive nor-
mal form 9l\ Vj=l~jVVj=]~;) is empty, the algo-
rithm is terminated. If not, then all conjunctions of this
disjunctive normal form are ordered in some way. For
the conjunction ~ which is the first in the sequence and
for all f from the set Fk(m), all disjunctive normal
forms which realize f on its domain of definition, which
consist of conjunctions in Sk(m, 9l) and which contain
the least number of variable symbols as compared with

other such disjunctive normal forms, are found.
Among them one selects all disjunctive normal forms
which, first, do not contain the conjunctions
~';, ... , ~~ and, secondly, contain all the conjunctions
~j, j= 1, ... ,s, for which Nw. n
(Ns \N)= 0. If for
all f in Fk(~ the conjunction ~ occurs in all disjunc-
tive normal forms selected, then the mark (/l, /l) over ~
is replaced by (1, /l) and the (i + l)-th step of the algo-
rithm is terminated. If ~ does not occur in any disjunc-
tive normal form selected, then the mark (/l, /l) is
changed into (Jl, 1) and the (i + l)-th step of the algo-
rithm is terminated. Otherwise, the procedure just
described is applied to the conjunction that is next in
sequence, etc. If the mark cannot be changed for any
conjunction, the algorithm terminates at the (i + l)-th
step. All conjun6tions obtained in the ring algorithm
over the abridged disjunctive normal form WI of a
function f with mark (1, /l) (or, respectively, (/l, 1
occur in all minimal disjunctive normal forms (or, do
not occur in any minimal disjunctive normal form) of f
The algorithms just described yield identical results,
whatever the manner of ordering of the conjunctions in
the disjunctive normal form.
The task of selecting all conjunctions occurring in at
least one or not occurring in any minimal disjunctive
normal form cannot be solved by algorithms working
with Sk(~, W) if k is bounded or increases too slowly as
the number of variables n increases. This situation does
not change if to the properties PI, P 2 stored in the
algorithm a bounded set of properties or one that
increases too slowly as n increases is added.
References
[I] YABLONSKII, S.V.: 'Functional constructions in k-placed logic',
Trudy Mat. Inst. Steklov. 51 (1958), 5-142 (in Russian).
[2] ZIflJRAVLEV, Yu.l.: 'Set-theoretical methods in the algebra of
logic', Problemy Kibemet. 8 (1962), 5-44 (in Russian).
[3] QUINE, W.V.: 'On cores and prime implicants of truth func-
tions', Amer. Math. Monthly 66 (1959),755-760.
Yu.I. Zhuravlev
AMS 1980 Subject Classification: 06E30
BOOLEAN FUNCI10NS, NORMAL FORMS OF -
Formulas of a special kind realizing Boolean functions.
One distinguishes between disjunctive normal forms (cf.
Boolean functions, minimization of) and conjunctive nor-
malforms. A product X'!l'I ... x''>'k ' where x" =x if CJ= 1
and x" = ex if CJ = 0, is called an elementary conjunction
of rank k if all its variables are different; '1' is con-
sidered to be an elementary conjunction of rank zero.
A logical sum xj: V ... V xj: is called an elementary
disjunction of rank r if all its variables are different; '0'
is considered to be an elementary disjunction of rank
zero.
A formula ~l V ... V~t. where ~J, , ~I are dis-
tinct elementary conjunctions of ranks rJ, ... , rl,
BOOLEAN FUNCTIONS, NORMAL FORMS OF
respectively, is called a disjunctive normal form, and
the number ~~ = 1ri is called its complexity; a formula
m1 ml> where m1, . , mt are different elementary
disjunctions of ranks PI, ... , Pt, respectively, is called
a conjunctive normal form, and the number ~:=IPi is
called its complexity. Every Boolean function that does
not vanish identically can be defined by a disjunctive
normal form which is, generally speaking, not unique.
The same applies to conjunctive normal forms and
functions that do not identically vanish.
From a table defining a Boolean function
f (x 1, . . . , x") one easily obtains the perfect disjunctive
normal form ~l V ... V~s, where ~i =X~il ... x~.. ,
i=l, ... ,s, and the sets CJil, ,CJi" are such that
f (CJi I, . . . , CJi") = 1. The perfect disjunctive normal form
realizing a Boolean function f is unique. The perfect
conjunctive normal form is defined similarly.
Since for 'almost-all' Boolean functions the number
of unit sets varies between 2" - I - Vn' 2" /2 and
2" -I + Vn' 2" /2, the asymptotic complexity of the per-
fect disjunctive normal form for 'almost-all' Boolean
functions is n2" -I. The maximal complexity of the
perfect disjunctive normal forms for functions in n vari-
ables is attained for functions that vanish at a single
point. This complexity is n '(2" - 1).
The main problem in the theory of normal forms of
Boolean functions is that of minimizing Boolean func-
tions, i.e. the construction of conjunctive or disjunctive
normal forms of minimal complexity - minimal con-
junctive or disjunctive normal forms for any given
Boolean function (cf. Boolean functions, minimization
of; Algorithm, local). Here, by virtue of the duality
principle, it is sufficient to consider disjunctive normal
forms only. In connection with the problem of minim-
izing Boolean functions one also considers abridged
dead-end, shortest and minimal disjunctive normal
forms. An important problem in the theory of disjunc-
tive normal forms is the search for numerical
characteristics of them, and characteristics connecting
various types of disjunctive normal forms of a given
function.
The abridged disjunctive normal form is uniquely
constructed from a Boolean function by means of a
fairly simple algorithm. Its most important property is
the fact that every minimal disjunctive normal form of
a function and at least one shortest form can be
obtained from the abridged disjunctive normal form by
the elimination of certain elementary conjunctions.
Therefore many minimization algorithms employ the
abridged disjunctive normal form as the initial specifi-
cation of a Boolean function. It is of major interest in
this context to define the complexity of abridged dis-
junctive normal forms for 'almost-all' functions and to
determine the absolute maximum of this complexity. If
429
BOOLEAN FUNCTIONS, NORMAL FORMS OF
Sin) is the number of elementary conjunctions in the
abridged disjunctive normal form of a Boolean function
!(xJ, ... ,xn ) and if
Sen) = max Sfin),
f(x" ... ,x,)
then the following estimates hold:
3n 3n
CI-::S S(n)::S C 2' _ I '
n vn
and for 'almost-all' Boolean functions
n(I-()log,log, n 2n < Sfin) < n(I+()log,log,n2n,
_0 as n_oo.
It is clear from these results and from the estimates of
the complexity of the perfect disjunctive normal form
that the complexity of the abridged disjunctive normal
form is much higher than that of the perfect disjunctive
normal form, both in the 'typical' and in the 'record-
breaking' cases. Unlike the perfect and the abridged
disjunctive normal form, a given Boolean function may
have many dead-end and minimal disjunctive normal
forms. Let tin) be the number of dead-end disjunctive
normal forms and let min) be the number of minimal
disjunctive normal forms of a Boolean function
!(xJ, ... ,xn ),
ten) = max tfin), men) =
f(x" . .. ,x,) f(x" ... ,x,)
Then the following estimates hold:
(22')Vn .;;; t(n)';;; (22'r/2, men) > (22')""nst.Vn,
and for 'almost-all' Boolean functions
(22'-' )(I-,)Iog, n log, log, n < tfin)
< (22H P+,)Iog, n log, log,
_0 as n_oo.
A non-trivial upper estimate for m(n) and a non-
trivial estimate for min) for 'almost-all' functions are
as yet (1977) not known. Complexity estimates of
dead-end and of minimal disjunctive normal forms are
of great interest in problems of minimization of
Boolean functions. Let II.J(n ) be the number of
elementary conjunctions in a dead-end disjunctive nor-
mal form T of the function !(x\, ... ,xn ); let lin) be
the number of elementary conjunctions in a shortest
disjunctive normal form of! (x \, ... ,xn ),
hen) = maxhJ(n),
f, T
len) = maxl;(n).
f
The following estimates then hold:
hen) ~ 2\ len) = 2"-1.
For 'almost-all' Boolean functions !(xJ, ... ,xn ) and
for almost-all dead-end disjunctive normal forms T:
hJ(n) ~ 2"-1.
For 'almost-all' Boolean functions! (x \, ... ,xn ):
430
2n 2n
lo~ n 'lo~ 10~ n < lfin) < 10~ n .
These estimates show that the shortest (and also the
minimal) disjunctive normal forms of 'almost-all'
Boolean functions account for only a small fraction of
the number of dead-end disjunctive normal forms.
There are also estimates of the relative complexity of
dead-end and shortest disjunctive normal forms of
Boolean functions. Let "hln) be the maximal number of
elementary conjunctions in a dead-end disjunctive nor-
mal form of a function !(xJ, ... ,xn ). For 'almost-all'
Boolean functions
10~ n < ~j=; < lo~ n '10~ 10~ n.
The maximum value of the ratio II.ln) / lin) can be
estimated from below by 2n (l-<), _0 as n_oo. Esti-
mates of y (f) - the so-called scatter of ! (x \, ... ,xn )
- have also been obtained. Here
y(j) = ~~ L(T")'
.MT.:.l
T,T
where T' and Tn are arbitrary dead-end disjunctive
normal forms realizing !(x\, ... ,xn ), and L(T) is the
number of characters in the dead-end disjunctive nor-
mal form T. Examples of Boolean functions with
max mfin). y (f)"~2n(\-o(\ have been constructed; it has been
established, however, that for 'almost-all' Boolean func-
tions
y(j) .;;; lo~ n 10& 10& n.
The estimates above give a fairly complete idea of the
difficulties that arise when Boolean functions are
< minimized according to the scheme: perfect disjunctive
n. normal form - abridged disjunctive normal form -
dead-end disjunctive normal form - minimal disjunc-
tive normal form.
References
[1] VASIL'EV, YU.L.: 'On comparing the complexity of typical and
disjunctive normal forms', Problemy Kibernet. 10 (1963),
5-61 (in Russian).
[2] GLAGOLEV, V.V.: 'Certain estimates of the disjunctive normal
forms of functions of the algebra of logic', Problemy Kibernet.
19 (1967), 75-94 (in Russian).
[3] KORSHUNOV, A.D.: 'The upper complexity bound of the shor-
test disjunctive normal forms of almost all Boolean functions',
Cybernetics 5, no. 6 (1969), 705-715. (Kibernetika (Kiev), no. 6
(1969), 1-8)
[4] SAPOZHENKO, A.A.: 'On the greatest length of a dead-end dis-
junctive normal form for almost all Boolean functions', Math.
Notes 4, no. 6 (1968), 881. (Mat. Zametki 4, no. 6 (1968),
649-658) Yu.l. Zhuravlev
AMS 1980 Subject Classification: 06E30
BOOLEAN RING - An associative ring K whose ele-
ments are all idempotent, i.e. x 2 =x for any x EK. Any
Boolean ring K=I=0 is commutative and is a subdirect
sum of fields ~ of two elements, and x + x = for all

x EK. A finite Boolean ring K=I=O is a direct sum of
fields ~ and therefore has a unit element.
A Boolean ring is the ring version of a Boolean alge-
bra, namely: Any Boolean algebra is a Boolean ring
with a unit element under the operations of addition
and multiplication defined by the rules
(x+y) = (xnCy)u(YnCx), xy =xny,
where Cx is the complement of x. The zero and the
unit of the ring are the same as, respectively, the zero
and the unit of the algebra. Conversely, every Boolean
ring with a unit element is a Boolean algebra under the
operations x Uy =x +y + xy, x ny =xy, Cx = 1+ x.
References
[I] STONE, M.H.: 'The theory of representations for Boolean alge-
bras', Trans. Amer. Math. Soc. 40 (1936),37-111.
[2] ZHEGALKIN, I.I.:"On the technique of computation of proposi-
tions in symbolic logic', Mat. Sb. 34, no. I (1927),9-28 (in
Russian). French abstract.
[3] VLADIMIROV, D.A.: Boolesche Algebren, Akademie Verlag, 1978
(translated from the Russian).
[4] SIKORSKI, R.: Boolean algebras, Springer, 1969.
Yu.M. Ryabukhin
Editorial comments. The operation
x+ y=(x n Cy) u
(y n Cx) is known as the symmetric
difference. Think of the Boolean algebra of all subsets of a
given set under union, intersection and complement to inter-
pret these formulas.
References
[A 1] RUDEANU, S.: Boolean functions and equations, North-
Holland,1974.
AMS 1980 Subject Classification: 06E20, 16A30,
16A32
BOOLEAN-VALUED MODEL - A model defined as
follows. Let ~ be the signature of some first-order
language L with one kind of variables, i.e. ~ is the set
of symbols of functions and predicates. A Boolean-
valued model then is a triple M=(BM' VM, ~M)' where
BM is a non-degenerate Boolean algebra, VM is a non-
empty set, and ~M is a function defined on ~ such that
0M(P) E V~M
if p is an n-place function symbol, and
0M(P) E BJ/
if p is an n-place predicate symbol. The symbol X Y
denotes the set of all functions defined on Y with
values in X and xn = x{m: m <n}, where n ~O is a
natural number. The Boolean algebra BM is called the
set of truth values of the model M. The set VM is called
the universe of M. A Boolean-valued model M is also
called a B-model if the set of truth values is the
Boolean algebra B, BM = B. If a Boolean algebra B is a
two-element algebra (i.e. B = {O, I}), then the B-model
M is the classical two-valued model.
Let LM be a language, complemented by new indivi-
BOOLEAN-VALUED MODEL
dual constants: each v E VM having its own individual
constant v. Let M be a B-model and let
B = (B; 0, 1, C, u, n) be a complete Boolean alge-
bra; the equalities 1) - 8) below then define the value
II e II M of each closed expression e (Le. of a formula or
a term without free variables) of L M :
1)11 v IIM=v, where VEVM ;
2)11 P(TJ, . ,tn ) IIM=(~M(p))(11 T\ 11M, ... , II Tn 11M),
where TJ, ,Tn are closed terms and p is an n-place
function or predicate symbol;
3) I/>:lljl IIM= -II I/> 11M U IlljI 11M;
4) I/>VljI IIM= II I/> 11M U IlljI 11M;
5) l/>/\ljI IIM= II I/> 11M n IIljI 11M;
6) -,1/> IIM= -II I/> 11M;
7) 3~q>(~) IIM= U VEV) q>(v) 11M;
8) V~q>(~ IIM= n II q>(v) 11M.
VEV
M
The relations 1) - 8) define the value II e II M for cer-
tain non-complete Boolean algebras as well; the only
condition is that the infinite unions and intersections in
7) and 8) exist. The concept of a Boolean-valued model
can also be introduced for languages with more than
one kind of variables. In such a case each kind of vari-
able has its own domain of variation VM.
A closed formula I/> is said to be true in a B-mode1 M
(MI=I/ if
A B-model M is said to be a model of a theory T if
MI=I/> for all axioms I/> of T. If h is a homomorphism of
a Boolean algebra B into a Boolean algebra B'
preserving infinite unions and intersections, then there
exists a B' model M' such that
for each closed formula I/> of L M If the universe of a
model M is countable, then there exists a homomor-
phism h into the Boolean algebra {O, I}, under which
M is transformed into the classical two-valued model
M' such that
It has been shown that a theory T is consistent if and
only if it has a Boolean-valued model. This theorem
forms the basis of the application of the theory of
Boolean-valued models to problems of the consistency
of axiomatic theories.
If the Boolean-valued model of a theory T is con-
structed by means of another axiomatic theory S, then
one obtains the statement on the consistency of T rela-
tive to S. Thus, the result due to P. Cohen on the con-
sistency of the theory ZF+(i~o>Nl) relative to ZF is
obtained by constructing the respective Boolean-valued
model by means of the system ZF (cf. Forcing method).
The construction of the Cohen forcing relation p I~I/> is
equivalent to that of a Boolean-valued model M such
431
BOOLEAN-VALUED MODEL

that while that bounded by a hyperbolic Booth lemniscate is

II cp II M = {P: p II- .-.c/>}.
a 2 -b 2 a ab
References S = 2 arctgb+T'
[1) RAsIOWA, E. and SIKORSKI, R.: The mathematics of
metamathematics, Polska Akad. Nauk, 1963. The Booth lemniscate is a special case of a Persian
[2) JECH, T.J.: Lectures in set theory: with particular emphasis on curve.
the method offorcing, Springer, 1971. Named after J. Booth [1).
[3) TAKE1ITl, G. and ZARING, W.M.: Axiomatic set theory,
Springer, 1973. References
[4) MANIN, Yu.I.: 'The problem of the continuum', J. Soviet [1) BOOTH, J.: A treatise on some new geometrical methods, 1-2,
Math. S, no. 4 (1976),451-502. (ftogi Nauk. i Tekhn. Sovrem. London, 1873-1877.
Problemy S (1975),5-73) [2) SAVELOV, A.A.: Plane curves, Moscow, 1960, pp. 144-146 (in
v.N. Grishin
Russian). D.D. Sokolov
Editorial comments.
AMS 1980 Subject Classification: 53A04
References
[A1) BELL, J.L.: Boolean-valued models and independence proofs
BORDERING ME1HOD - A method for solving a
in set theory, Clarendon Press, 1977.
[A2) JECH, T.J.: Set theory, Acad. Press, 1978. system of linear algebraic equations Ax =b with a non-
[A3) KUNEN, K.: Set theory, North-Holland, 1980. degenerate matrix by inverting a matrix and calculating
AMS 1980 Subject Classification: 03C90 a determinant, based on passing recursively from the
solution of the problem with matrix
BOOTH LEMNISCATE - A plane algebraic curve of
order four whose equation in orthogonal Cartesian
coordinates is
ak-I.I ak-I.k-I
(x 2 +y2)2 +(2m 2 +n)x 2 +(2m 2 -n)y2 = o.
to the solution of the problem with the matrix A k ,
If I n I <2m 2, the Booth lemniscate is called elliptic (it
which can be considered as the result of bordering
has singular point 0 (Fig. 1), where O<n<2m2). If
Ak -
I n I > 2m 2 , the Booth lemniscate is called hyperbolic (it I
The calculating scheme of the bordering method for
has a nodal point at the coordinate origin, cf. Fig. 2,
the inversion of a matrix is as follows: Let Ak -I be a
where n >2m 2 ).
non-degenerate matrix. In the inversion of the matrix
y
A k , the representation

IIAk-1 I

ob
Ak = Vk
Uk
akk (1)
o x
is used, where Uk =(a I,k> ... ,ak + I,kf and
..x Vk=(ak , I, . . . ,akk-I).
, Then
A-I A-I
-I k-IUkVk k-I Ak~IUk
A k- I +
Ci.k Ci.k
Fig. l. Fig. 2. Ak l (2)
vkAk~1
Ci.k Ci.k
The equation of an elliptic Booth lemniscate in polar
coordinates is Ci.k = akk-vkAk~IUk'
p2 = a 2 cos2 cp+ b 2 sin2 cp or p o. Subsequent inversion of the matrices A I, . . . ,An by
this scheme gives the matrix A -I .
If n > 2m 2, the equation of a hyperbolic Booth lemnis-
The above scheme of the bordering method is only
cate has the form
applicable to matrices with principal minors unequal to
p2 = a2cos2cj>-b2sin2cj>; zero. Generally, a choice of the principal element is
and if n < - 2m 2 adopted. In this scheme, the bordering rows and
p2 = -a 2 cos 2 cj>+b 2 sin2 cj> columns used are those for which ak=akk-vkAk~IUk
has the maximum absolute value. The calculated matrix
(a 2 = 12m 2 +n I, b2 = 12m 2 -n I). will then differ from A -I only in the rearrangement of
The arc length of a Booth lemniscate is expressed by the rows and columns [1].
elliptic integrals. The area bounded by an elliptic The hordering method is not the fastest of the direct
Booth lemniscate is methods\ of inversion of a matrix.
The bordering method permits an effective inversion
of a triangular matrix. If Ak is a right (upper-) triangu-

432
 BORDERING METHOD

lar matrix, then in (1) ing method is equivalent to the Gauss method, one of
the fastest direct methods of solving systems.
The bordering method allows one to solve higher-
order systems owing to the effective use of computer
o memory. This is caused by the fact that for the calcula-
tion of the vectors x (k,p) , P > k, only the storage of the
The amount of calculation in this instance is reduced vectors x(k -I,p), P >k -1, and the coefficients of a k-th
six-fold. order system of equations is necessary, that is, a series
The bordering method is particularly effective when of numbers of length j(k)=k(n - k + l)+(n + 1). It is
inverting Hermitian positive-definite matrices. For these therefore sufficient to have a working capacity
matrices it is not necessary to use a scheme with a (n + 1Xn + 5) /4::::::(n /2)2 to solve a system of order n.
choice of a principal element. Moreover, they are Moreover, the elements of the matrix and of the right-
determined by only half of their elements. The calcu- hand side do not have to be put into the computer
lating scheme in this instance is simplified thus: memory immediately but successively, row by row.
A-I
pr
+-k k
Pk It is advisable to use the bordering method when
k-I -
ak ak solving systems for which a truncated system has
Ak l already been solved. Relation (4) then immediately
Pi,
ak ak
gives the required solution.
The scheme of the bordering metho.d described can
be used to calculate the determinant. It follows from
The calculating scheme of the bordering method for (1) that
the solution of a system is as follows. Let
1 Ak 1 = (akk-vkAk~IUk) 1 Ak- I I,

b (kp)-(
- alp' ... ,akp')T k=l , ... , n', b(n,n+I)=-b . If Recurrent use of this relation gives I A I
Ak -I is a non-degenerate matrix and x(k -I,p) is the As with the inversion of a matrix, the solution of a
solution of the system Ak_IX(k-l,p)+b(k-l,p)=O, then system and the calculation of the determinant by means
the solution X(k,p) of the system Akx(k,p) +b(k,p) =0 is of the bordering method is possible only for matrices
obtained from the representation with non-zero principal minors. It is generally also
necessary here to use a scheme with a choice of a prin-
Ak-1 b(k-I,k) II = Ilb(k-I,P) I
A
k
=
II Vk akk '
b(k,p)
akp ,
(3) cipal element.
References
and from (2) as follows [1) VOEVODIN, V,V.: Numerical methods of algebra, Moscow, 1966

x (k,p) = -
I Ak~1 llllb(k-I,P) II + (in Russian).
(2) FADDEEv, D.K. and FADDEEVA, V,N,: Computational methods
o 0 akp
of linear algebra, Freeman, 1963 (translated from the Russian).
G.D. Kim
Editorial comments. The bordering method is a special
case of more general partitioning procedures [A 1].
= Ilx(k-I,P) 11_ VkX(k-l,p)+akp Ilx(k-l,k) II
The bordering method defined above is related to the
o akk+vkx(k-I,k) 1 (4) Sherman- Morrison formula and its generalization the
Sherman- Morrison- Woodbury formula. The
In this way, through the solutions x(k -I,p) and x(k -I,k) Sherman - Morrison formula states that a non-singular
of the systems with the same matrix Ak -I and with (m X m)-matrix M and vectors x and yare related by
various right-hand sides, it is easy to obtain the solu- [M+xyHr' = M-'-M-'xc-'yHM-', C=1+yHM-'x,
tion of the system with the bordered matrix A k The
provided that C=FO. The Sherman - Morrison - Woodbury
solution of the initial sytem is x(n,n + I). This can be
formula replaces the vectors x and y by (mxn)-matrices X
obtained by the recurrent use of relation (4). This leads and Y. These formulas were developed in [A2] - [A5] (see
to the subsequent calculation of the set of vectors x (k,p) , also [A6]) , and used for the construction of matrix inversion
k= 1, ... ,n,p>k, that is schemes resembling that of the bordering method. An
X(12), X(l3), X(I,n+l) , excellent discussion of these schemes is given in [A7].

X(23), x(2,n + I), References

[A1] WESTLAKE, l.R.: A handbook of numerical matrix inversion
and solution of linear equations, Wiley, 1968, Sect. 2.6.
x(n -I,n +I), [A2] BARTLETT, M.S.: 'An inverse matrix adjustment arising in
x(n,n+I), discriminant analysis', Ann. Math. Stat. 22 (1951),107-111.
[A3] SHERMAN, l.: Computations relating to inverse matrices, Appl.
In the amount of calculation necessary, this border- Math., 29, Nat. Bur. Standards, 1953, pp. 123-124.

433
BORDERING METHOD
[A4] SHERMAN, J. and MORRISON, W.J.: 'Adjustments of an inverse
matrix corresponding to changes in the elements of a given
column or a given row of the original matrix', Ann. Math.
Stat. 20 (1949), 621.
[A5] WOODBURY, M.: 'Inverting modified matrices', Memorandum
Report Stat. Res. Group Princeton 42 (1950).
[A6] GoLUB, G.H. and loAN, C.F. VAN: Matrix computations,
North Oxford Acad., 1983.
[A7l HOUSEHOWER, A.S.: Principles of numerical analysis,
McGraw-Hili, 1953.
AMS 1980 Subject Classification: 65F05, 65F10,
15A09
BORDERING OF A SPACE X in a compactification
bX - A finite family {U 1, . . . , Uk} of sets open in X
such that the set K~X\ (U I U .~.. U Uk) Js com-
pact, and bX=K U U I U ... U Ub where Uj is the
largest open set in bX the intersection of which with X
is the set U; (X is assumed to be completely regular).
The concept of a bordering of a space X in bX coin-
cides with the concept of an almost-extendable border-
ing of a proximity space X (the proximity on X is
induced by the extension bX), formulated in terms of
the proximity: apart from K being compact, it is neces-
sary that for any neighbourhood (f) K, the family
{ (f) K, U J, , Uk} is a uniform covering of the space
X. A bordering of a space X in its Stone - Cech com-
pactification is simply called a bordering of X. In the
language of borderings, a series of theorems has been
formulated on the dimensions of the remainder of com-
pactifications of topological and proximity spaces.
References
[II SMIRNOV, Yu.M.: 'On the dimensions of remainders of com-
pactifications of proximity and topological spaces', Mat. Sb.
71, no. 4 (1966), 554-482 (in Russiap).
v. V. Fedorchuk
Editorial comments. A concept related to the bordering
of a space is that of a border cover. A collection <t of open
sets such that X\ U <t is compact. Border covers work in
a sense opposite to borderings. In the case of borderings a
compactification is given; from certain systems of border
covers one can construct compactifications whose
remainders can have special properties.
AMS 1980 Subject Classification: 54020, 54035,
54040
BORDISM, bordantism - A term used by itself or as
a part of standard expressions in a number of similar
meanings. The older term, cobordism, is still employed.
The simplest variant is the following. Two smooth
closed n-dimensional manifolds M 0 and M are bordant
(cobounding or internally homologous) if there exists a
smooth compact (n + I)-dimensional manifold W (a
'membrane') the boundary of which consists of two
manifolds No and N which are diffeomorphic, respec-
tively, to M 0 and M under certain diffeomorphisms (cf.
Diffeomorphism)
434
go:No ~ Mo and g:N ~ M.
A set of mutually bordant manifolds is called a bordism
class, while the triplet (W, M o, M) is sometimes called
a bordism (it would be more accurate to take
(W, Mo, M, go, g. The set of bordism classes of n-
dimensional manifolds forms an Abelian group ~" with
respect to unconnected union. The zero element of the
group is the bordism class consisting of the manifolds
M which constitute the boundary of a certain manifold
W (one formally speaks of the triplet (W, M o, M) with
empty M 0; other appellations are: M is a 'bounding
manifold', or M is 'internally homologous to zero' or is
'bordantly zero'). The element of ~" inverse to a given
bordism class is this class itself (since MUM is dif-
feomorphic to the boundary of the direct product
M X [0, I D. The direct sum ~. of the groups ~" is a
commutative graded ring in which the multiplication is
induced by the direct product of manifolds with unit
element given by the bordism class of a point.
More complex variants comprise the bordism of
smooth closed manifolds with a supplementary struc-
ture. For instance, two oriented manifolds M 0 and M
are orientedly bordant if they are bordant in the sense
explained above, if the 'membrane' W is oriented and
if, using the above notation, the orientation induced by
the orientation of W on No and N (which are parts of
the boundary) goes over under the diffeomorphisms go
and g, respectively, to the original orientation on M
and the opposite orientation on Mo. One then speaks
of oriented bordism; if it is desired to stress the differ-
ence between such bordism and bordism in the sense
given above, then one speaks of the latter as non-
oriented bordism. In analogy with ~n and ~. one intro-
duces the groups of oriented bordism nnand the ring
n.=~n".
Historically, the first example was that of the bor-
dism of nested manifolds, introduced in 1938 by L.S.
Pontryagin, who showed that the classification of such
bordisms is equivalent to the calculation of the homo-
topy groups of the spheres 7Ti(S"), and who could
determine 7T" + 1(S") and 7T" +2(S") in this way (see [2]
for a detailed account of his studies, and [4] for an
introductory text). Non-oriented and oriented bordisms
were introduced in 1951 - 1953 by V.A. Rokhlin [3],
who calculated ~" and nn for n ~4. It had been previ-
ously shown by Pontryagin [1] that if two manifolds are
bordant, their characteristic numbers are identical
(Stiefel- Whitney numbers for non-oriented manifolds,
Stiefel- Whitney and Pontryagin numbers for oriented
manifolds). It was subsequently found that the converse
proposition is also true.
Modem methods of algebraic topology were first
applied in the theory of bordism in 1954 by R. Thorn
 BOREL - CANTELLI LEMMA

[5], [6], who rediscovered (for the case of both oriented 9Cn(X) and each mapping I> with the mappings I>n is a
and non-oriented bordisms) the connection between generalized homology theory. In the present case it
bordisms and certain homotopy problems. Thus, the reduces to ordinary homology, viz., for any cellular
group 9Cn is isomorphic to the group '!Tn +r(TBO(r)) for polyhedron X
a sufficiently large r; here, TBO(r) is the Thorn space of L9?n(X) ~ H.(X, ~)9?
the universal vector bundle with structure group OCr).
Owing to this connection, Thorn was able to make a (the tensor product of graded modules on the right-
complete computation of the ring 9C. and to contribute hand side of this expression is over Z:J. = Z / (2)), but
substantially to the study of ~" which was subse- this is generally not true for other bordisms (oriented,
quently continued by other workers. Thus, 9C. proved etc.). Many generalized homology theories may be
to be the ring of polynomials over the field of residues obtained by way of a so-called bordism with singulari-
modulo 2 in the generators Xi of dimension i, where i ties [14].
runs through all positive numbers not equal to 2S -1, In addition to the bordism of a space X there exist
s ~ 1; a geometric realization of these generators is generalized cohomology theories which are dual to it.
known (i.e. concrete manifolds whose bordism classes The introduction of these generalized homology and
are Xi have been given [7]). cohomology functors has made it desirable to introduce
Other variants of bordism of manifolds with a sup- the changed terminology discussed at the beginning of
plementary structure comprise the very important bor- the present article. Thus, the term 'cobordism' is res-
disms of quasi-complex manifolds (also called unitary tricted to generalized cohomology theories dual to bor-
bordism or complex cobordism [8], [9]), and bordism of dism.
manifolds acted upon by a group of transformations References
[10]. There are also variations of another kind (for [1) PONTRYAGIN, L.S.: 'Characteristic cycles of differentiable man-
ifolds', Mat. Sb. 21 (63), no. 2 (1947), 233-284 (in Russian).
piecewise-linear or topological manifolds, for Poincare [2) PONTRYAGIN, L.S.: Smooth manifolds and their applicatiOns in
complexes, etc. [11]). Special kinds of bordisms include homology theory, Moscow, 1976 (in Russian).
foliated bordism and h-bordism (previously referred to as [3) ROKHLIN, V.A.: 'Theory of inner homology', Uspekhi Mat.
Nauk 14, no. 4 (1959), 3-20 (in Russian).
J-equivalences); these are used for relating differential [4A) WALLACE, A.H.: Differentiable topology. First steps, Benjamin,
and homotopy topological properties [12]. 1968.
Further development of bordism theory is related to [4B) MILNOR, J.W.: Topology from the differentiable viewpoint, Uniy.
Press of Virginia, 1965.
the bordism groups of a topological space X (bordism of [5) THOM, R.: 'Quelques proprietes globales des yarietes
spaces for short, cf. [13]). They are defined for dif- differentiables', Comm Math. Helvetia 28 (1954), 17-86.
ferent variants of bordisms (the simplest example is [6) MILNOR, J.W. and STASFEFF, J.D.: Characteristic classes,
given below). A singular n-dimensional (sub)manifold of Princeton Uniy. Press & Uniy. of Tokyo Press, 1974.
[7) DoLO, A.: 'Erzeugende der Thomschen Algebra 9?',
a space X is a pair (Mn ,f), where M n is a closed Math. Z. 65 (1956),25-35.
smooth manifold, and f: Mn~x is a continuous map- [8) MILNOR, J.: 'On the cobordism ring (r and a complex analo-
ping. Two such pairs (M o,fo), (Mn,f) are bordant if gue 1', Amer. J. Math. 82 (1960),505-521.
[9) NOVDWV, S.P.: 'Homotopic properties of Thorn-complexes',
M 0 and M are bordant in the ordinary sense, and Mat. Sb. 57 (99) (1962), 406-442 (in Russian).
(using previous notations) if there exists a continuous [10) CONNER, P.E. and FLOYD, E.E.: Differentiable periodic maps,
o'
mapping h: W~X such that hg =fo, hg-' =fo. (If Springer, 1964.
[11) STONG, R.E.: Notes on cobordism theory, Princeton Uniy. Press,
M 0 and M are identified with No and N, one can sim- 1968.
ply say that the mapping h induces the given mappings [12) MILNOR, J.: Lectures on the h-cobordism theorem, Princeton
of M 0 and Minto x.) The bordism classes of singular Uniy. Press, 1965.
manifolds in a space X form the n-dimensional bordism [13) ATIYAH, M.F.: 'Bordism and cobordism', Proc. Cambridge Phi-
los. Soc. (2) 57 (1961), 200-208.
group 9Cn (X) of this space (the group operation is gen- [14A) BAAS, N.A.: 'On bordism theory of manifolds with singulari-
erated by the union of manifolds). ties', Math. Scand. 33 (1973), 279-302.
If X is a manifold of dimension > 2n, the elements of [14B) BAAS, N.A.: 'On formal groups and singularities in complex
bordism theory', Math. Scand. 33 (1973),302-313.
9Cn (X) can be visualized as submanifolds, as can the
D. V. Anosov
corresponding 'membranes'; in this respect bordism of
M.I. Voitsekhovskii
spaces resembles the original attempts at introducing
homologies. If X is a point, 9Cn (X) is reduced to the AMS 1980 Subject Classification: 57N70
previous 9Cn. To a mapping 1>: X~y correspond
homomorphisms I>n: 9Cn(X)~9Cn(Y)' generated by the BOREL - CANTELLI LEMMA - A frequently used
transition from a singular manifold (M, f) to another statement on infinite sequences of random events. Let
one (M, I>f). The functor by means of which each A" ... ,An, ... , be a sequence of events from a cer-
space X is brought into correspondence with the groups tain probability space and let A be the event consisting

435
BOREL - CANTELLI LEMMA
in the occurance of (only) a finite number out of the
events An, n =1, 2, . ... Then, according to the
Borel- Cantelli lemma, if
(*)
then
P(A) = 1.
If the events An are mutually independent, then
P(A) = 1 or 0, depending on whether the series
~:=l P(An) converges or diverges, i.e. in this case the
condition (*) is necessary and sufficient for P(A) = 1;
this is the so-called Borel criterion for 'zero or one' (cf.
Zero-one law). This last criterion can be generalized to
include certain classes of dependent events. The
Borel- Cantelli lemma is used, for example, to prove
the strong law of large numbers.
References
[I) BoREL, E.: 'Les probabilites denombrables et leurs applications
arithmetique', Rend. Circ. Mat. Palermo (2) 27 (1909),247-271.
[2) CANfELLI, F.P.: 'Sulla probabilita come limite della frequenza',
Atti Accad. Naz. Lincei 26, no. I (1917), 39-45.
[3) LoEVE, M.: Probability theory, Princeton Univ. Press, 1963.
A. V. Prokhorov
Editorial comments. The Borel- Cantelli lemma can be
used in number theory to prove the so-called 'normality' of
almost-all natural numbers, cf. [A1], Chapt. 8, Sect. 6.
References
[A1) FELLER, W.: An introduction to probability theory and its
applications, 1, Wiley, 1957.
AMS 1980 Subject Classification: 60A05
BOREL FIELD OF EVENTS, a-field, Borel algebra,
a-algebra of events - A class A of subsets (events) of a
non-empty set n (the space of elementary events) which
is a Borel field of sets.
v. V. Sazonov
AMS 1980 Subject Classification: 28A05
BOREL FIELD OF SETS, family of Borel sets, gen-
erated by a system of sets M - The smallest system of
sets containing M and closed with respect to the opera-
tions of countable union and taking complements.
A.G. El'kin
AMS 1980 Subject Classification: 54C50, 54C99
BOREL FIXED-POINT THEOREM - A connected
solvable algebraic group G acting regularly (cf. Alge-
braic group of transfonnations) on a non-empty com-
plete algebraic variety V over an algebraically-closed
field k has a fixed point in V. It follows from this
theorem that Borel subgroups (cf. Borel subgroup) of
algebraic groups are conjugate (The Borel- Morozov
theorem). The theorem was demonstrated by A. Borel
[l]. Borel's theorem can be generalized to an arbitrary
436
(not necessarily algebraically-closed) field k: Let V be a
complete variety defined over a field k on which a con-
nected solvable k-split group G acts regularly, then the
set of rational k-points V(k) is either empty or it con-
tains a point which is fixed with respect to G. Hence
the generalization of the theorem of conjugation of
Borel subgroup is: If the field k is perfect, the maximal
connected solvable k-split subgroups of a connected k-
defined algebraic group H are mutually conjugate by
elements of the group of k-points of H [2].
References
[1) BOREL, A.: 'Groupes lineaires algebriques', Ann. of Math. (2)
64, no. 1 (1956), 20-82.
[2) BOREL, A.: Linear algebraic groups, Benjamin, 1969.
[3] MOROZOV, V.V.: Dokl. Akad. Nauk SSSR 36, no. 3 (1942),
91-94.
v.P. Platonov
AMS 1980 Subject Classification: 22A05
BOREL FUNcrION, B-function - A function for
which all subsets of the type E(x: f(x)~a) in its
domain of definition are Borel sets (cf. Borel set). Such
functions are also known as Borel-measurable functions,
or B-measurable functions. The operations of addition,
multiplication and limit transition - as in the general
case of measurable functions - do not take one out-
side the class of Borel functions, but, unlike in the gen-
eral case, the superposition of two Borel functions does
also not lead outside the class of Borel functions. More-
over, [1], if f is a measurable function on an arbitrary
space n and if g is a Borel function on the space of real
numbers, then the function h =g[f(x)] is measurable
on n. All Borel functions are Lebesgue-measurable (cf.
Measurable function). The converse proposition is not
true. However, for any Lebesgue-measurable function f
there exists a Borel function g such that f(x)=g(x)
almost-everywhere [1]. Borel functions are sometimes
called Baire functions, since the set of all Borel func-
tions is identical with the set of functions belonging to
the Baire classes (Lebesgue's theorem, [2]). Borel func-
tions can be classified by the order of the Borel sets
E(x: f(x)~a); the classes thus obtained are identical
with the Baire classes.
The concept of a Borel function has been generalized
to include functions with values in an arbitrary metric
space [3]. One then also speaks of B-measurable map-
pings. Borel functions have found use not only in set
theory and function theory but also in probability
theory [1], [4].
References
[I) HALMos, P.R.: Measure theory, v. Nostrand, 1950.
[2] HAUSDORFF, F.: Set theory, Chelsea, reprint, 1978.
[3) KURATOWSKI, K.: Topology, Acad. Press, 1966-1968 (translated
from the French).
[4) KOIMOGOROV, A.N.: Foundations of the theory ofprobability,

Chelsea, reprint, 1950 (translated from the Russian).
VA. Skvortsov
AMS 1980 Subject Classification: 26A21, 26A30,
28A20, 28A05
BOREL ISOMORPHISM, B-isomorphism - A one-
to-one mapping f of a space X into a space Y such that
both f and f- 1 transform Borel sets into Borel sets (cf.
Borel set). In the class of Borel subsets of complete
separable metric spaces, sets of the same cardinality are
Borel isomorphic.
A.G. El'kin
AMS 1980 Subject Classification: 54C50, 54C99
BOREL - LEBESGUE COVERING THEOREM - Let
A be a bounded closed set in R n and let G be an open
covering of it, i.e. a system of open sets the union of
which contains A; then there exists a finite subsystem
of sets {G;}, i = 1, ... ,N, in G (a sub covering) which
is also a covering of A, i.e.
N
A C UG,.
i=!
The Borel- Lebesgue theorem has a converse: If
A eRn and if a finite sub covering may be extracted
from any open covering of A, then A is closed and
bounded. The possibility of extracting a finite subcov-
ering out of any open covering of a set A is often taken
to be the definition of the set A to be compact. Accord-
ing to such a terminology, the Borel- Lebesgue
theorem and the converse theorem assume the follow-
ing form: For a set A eRn to be compact it is neces-
sary and sufficient for A to be bounded and closed.
The theorem was proved in 1898 by E. Borel [1] for the
case when A is a segment [a, b] e R 1 and G is a system
of intervals; the theorem was given its ultimate form by
H. Lebesgue [2] in 1900 - 1910. Alternative names for
the theorem are Borel lemma, Heine-Borel lemma,
Heine- Borel theorem.
References
[I) BOREL, E.: Ler;ons sur la theorie des fonctions, Gauthier-Villars,
1928.
[2) RUDIN, W.: Principles of mathematical analySiS, McGraw-HilI,
1953.
I.A. Vinogradova
AMS 1980 Subject Classification: 54030
BOREL MEASURE of sets - A non-negative function
J.L of the subsets of a topological space X possessing the
following properties: 1) its domain of definition is the
a-algebra !J# of Borel sets (cf. Borel set) in X, i.e. the
smallest class of subsets in X containing the open sets
and closed with respect to the set-theoretic operations
performed a countable number of times; and 2)
J.L( U ~=lEi)="2:~=lJ.L(Ei) if E; nEj= 0 when i-=Fj, i.e.
BOREL MEASURE
is countably additive. A Borel measure
J.L J.L is called reg-
ular if
11.(.E) = sup- p.(F),
E:JF
where F belongs to the class $' of closed subsets in X.
The study of Borel measures is often connected with
that of Baire measures, which differ from Borel meas-
ures only in their domain of definition: They are
defined on the smallest a-algebra !J# 0 with respect to
which all continuous functions on X are measurable. A
Borel measure J.L (or a Baire measure v) is said to be T-
smooth if J.L(Fa)tO for any net {Fa} of closed sets which
satisfies the condition Fat 0 (or v(ZaHO for any net
{Z a} of sets which are zero sets of continuous func-
tions and such that Z at 0 ). A Borel measure J.L (or
Baire measure v) is said to be tight if
p.(E) = E :JKEf
sUP p.(K),
where .Y( is the class of compact subsets on X (or
P(E) = E :JKEf
sUP p' (K),
where
Tightness and T-smoothness are restrictions which
ensure additional smoothness of measures, and which
in fact often hold. Under certain conditions Baire
measures can be extended to Borel measures. For
instance, if X is a completely-regular Hausdorff space,
then any T-smooth (tight) finite Baire measure can be
extended to a regular T-smooth (tight) finite Borel
measure. In the study of measures on locally compact
spaces Borel measures (or Baire measures) is the name
sometimes given to measures defined on the a-ring of
sets generated by the compact (or G8 -compact) sets
and which are finite on compact sets. Often, by the
Borel measure on the real line one understands the
measure defined on the Borel sets such that its value on
an arbitrary segment is equal to the length of that seg-
ment.
References
[1] VARADARAJAN, V.S.: 'Measures on topological spaces', Transl.
Amer. Math. Soc. Ser. 248 (1965), 161-228. (Mat. Sb. 55 (97),
no. 1 (1961),35-100)
[2] HALMOS, P.R.: Measure theory, v. Nostrand, 1950.
[3] NEVEU, J.: Bases mathematiques du calcul des probabi/ites, Mas-
son, 1970.
VV Sazonov
Editorial comments.
References
[A1] ROYDEN, H.L.: Real analysis, Macmillan, 1968.
[A2] ZAANEN, A.C.: Integration, North-Holland, 1967.
[A3] RUDIN, W.: Principles of mathematical analysis, McGraw-Hili,
1953.
[A4] RUDIN, W.: Real and complex analysis, McGraw-Hili, 1966.
[AS] TAYLOR, A.E.: General theory of functions and integration,
Blaisdell, 1965.
437
BOREL MEASURE
[A6] ALIPRANTZ, C.D. and BURLEINSHAW, 0.: Principles of real
analysis, North-Holland, 1981.
AMS 1980 Subject Classification: 28A12
BOREL SET, B-set - A set which may be obtained
as the result of not more than a countable number of
operations of union and intersection of closed and open
sets in a topological space. More exactly, a Borel set is
an element of the smallest countably-additive class of
sets containing the closed sets, and which is closed with
respect to complementation. Borel sets are also called
Borel-measurable sets and B-measurable sets. Open and
closed sets are said to be Borel sets of order zero. Borel
sets of order one are sets of type F" or G 8 which are,
respectively, countable sums of closed sets and count-
able intersections of open sets. Borel sets of the second
order are sets of type F ,,8 (the intersection of a count-
able number of sets of type F (J) and sets of type Goo
(the union of a countable number of sets of type G8 ).
Borel sets of an arbitrary finite order are defined in a
similar manner by induction. With the aid of transfinite
numbers of the second class (cf. Transfinite number)
this classification may be exhaustively extended to all
Borel sets. Let a be an arbitrary transfinite number of
the second class; the Borel sets of class a will include
any Borel set of order a that is not a Borel set of order
a' for any a' <a. Whether the classes of Borel sets are
empty or not will depend on the basic space under con-
sideration. In Euclidean, Hilbert and Baire spaces there
exist Borel sets of all classes.
Borel sets are a special case of d -sets. For an d -set
E to be a Borel set it is necessary and sufficient that
the complement of E also is an d -set (Suslin's cri-
terion). In spaces with a Lebesgue measure all Borel
sets are Lebesgue measurable. The converse proposition
is not true. All separable spaces having the cardinality
of the continuum contain sets that are not Borel sets.
Borel sets were introduced by E. Borel [1]; they play
an important role in the study of Borel functions (cf.
Borel function).
In a more general sense a Borel set is a set in an
arbitrary Borel system of sets generated by some sys-
tem of sets. The Borel sets in a topological space are
generated by the system of closed subsets of this space.
References
[I] BOREL, E.: Lefons sur les fonctions discontinues, Gauthier-
Villars, 1898.
[2] KURATOWSKI, K.: Topology, Acad. Press, 1966-1968 (translated
from the French).
[3] HAUSDORFF, F.: Set theory, Chelsea, reprint, 1978.
[4] ALEKSANDROV, P.S.: EinfiIhrung in die Mengenlehre und die
Theorie der reelen Funktionen, Deutsch. Verlag Wissenschaft.,
1956 (translated from the Russian).
V.A. Skvortsov
Editorial comments. For notational questions, see Borel
set of ambiguous class.
438
References
[A1] KURATOWSKI, K.: Introduction to set theory and topology,
Pergamon, 1961.
AMS 1980 Subject Classification: 54C50, 54C99
BOREL SET, CRITERION FOR A - A necessary and
sufficient condition for an d -set in a complete separ-
able metric space to be a Borel set. Two criteria can be
stated as follows: 1) its complement must also be an
d-set (Suslin's criterion); and 2) it can be represented
as the union of non-intersecting components (Luzin's
criterion).
A.G. El'kin
AMS 1980 Subject Classification: 54C50
BoREL SET OF AMBIGUOUS CLASS a - A Borel
subset of a metric, or (more general) of a perfectly-
normal topological, space that is at the same time a set
of additive class a and of multiplicative class a, i.e.
belongs to the classes Fa and Ga at the same time.
Borel sets of ambiguous class 0 are the closed and open
sets. Borel sets of ambiguous class I are sets of types
F" and G8 at the same time. Any Borel set of class a is
a Borel set of ambiguous class fi for any {3>a. The
Borel sets of ambiguous class a form a field of sets.
References
[1) KURATOWSKI, K.: Topology, 1, Acad. Press, 1966 (translated
from the French).
[2) HAUSDORFF, F.: Set theory, Chelsea, reprint, 1978 (translated
from the German). A. G. El'kin
Editorial comments. The notations Fa, Ga are still
current in topology. Outside topology one more often uses
the notation ~2, II2, respectively. For a~w one has
Fa=~2, Ga =II2; but ~g+1=Fn and IIg+ 1 =Gn for n<w.
The notation for the ambiguous classes is ~2 =~2 II~. n
See also [A 1].
References
[A1] MOSCHOVAKIS, Y.: Descriptive set theory, North-Holland,
1980.
AMS 1980 Subject Classification: 54H05, 04A15
BOREL STRONG LAW OF LARGE NUMBERS -
Historically, the first variant of the strong law of large
numbers, formulated and proved by E. Borel [1] in the
context of the Bernoulli scheme (cf. Bernoulli trials).
Consider independent random variables
Xi> ... ,Xn , . . . , which are identically distributed and
assume one of two values 0 and 1 with probability of
112 each; the expression Sn =~~ = I Xk will then give
the number of successful trials in a Bernoulli scheme in
which the probability of success is 112. Borel [1]
showed that
Sn 1
-~
n 2
with probability one as n~oo. It was subsequently

(1914) shown by G.H. Hardy and J.E. Littlewood that,
almost certainly,
lim
n~oo
Is"-~
yin
I<
V2'
Inn
_1_
l=l
after which (1922) the stronger result:
-!!...
n 2
P -li
m - -I-
- VP. Platonov
n-->oo Yn In Inn V2
was proved by A.Ya. Khinchin. See also Law of the
iterated logarithm.
References
[I] BOREL, E.: 'Les probabilites denombrabies et ieurs applications
arithmetique', Rend. Circ. Mat. Palermo (2) 27 (1909),247-271.
[2] KAc, M.: Statistical independence in probability, analysis and
number theory, Math. Assoc. Amer., i963. A V P kh
. . ro orov
AMS 1980 Subject Classification: 60F15
BOREL SUBGROUP - A maximal connected solv-
able algebraic subgroup of a linear algebraic group G.
Thus, for instance, the subgroup of all non-singular
upper-triangular matrices is a Borel subgroup in the
general linear group GL(n). A. Borel [I] was the first to
carry out a systematic study of maximal connected
solvable subgroups of algebraic groups. Borel sub-
groups can be characterized as minimal parabolic sub-
groups, i.e. algebraic subgroups H of the group G for
which the quotient variety G / H is projective. All
Borel subgroups of G are conjugate and, if the Borel
subgroups B I, B 2 and the group G are defined over a
field k, B I and B 2 are conjugate by an element of
G(k). The intersection of any two Borel subgroups of a
group G contains a maximal torus of G; if this intersec-
tion is a maximal torus, such Borel subgroups are said
to be opposite. Opposite Borel subgroups exist in G if
and only if G is a reductive group. If G is connected, it
is the union of all its Borel subgroups, and any para-
bolic subgroup coincides with its normalizer in G. In
such a case a Borel subgroup is maximal among all
(and not only algebraic and connected) solvable sub-
groups of G. Nevertheless, maximal solvable subgroups
in G which are not Borel subgroups usually exist. The
commutor subgroup of a Borel subgroup B coincides
with its unipotent part B u, while the normalizer of Bu
in G coincides with B. If the characteristic of the
ground field is 0, and g is the Lie algebra of G, then
the subalgebra b of g which is the Lie algebra of the
Borel subgroup B of G is often referred to as a Borel
subalgebra in g. The Borel subalgebras in g are its
maximal solvable subalgebras. If G is defined over an
arbitrary field k, the parabolic subgroups which are
defined over k and are minimal for this property, play
BOREL SUMMA nON METHOD
a role in the theory of algebraic groups over k similar
to that of the Borel groups. For example, two such par-
abolic subgroups are conjugate by an element of G(k)
[2].
References
[i] BOREL, A.: 'Groupes lineaires algebriques', Ann. of Math. (2)
64, no. I (1956), 20-82.
[2] BOREL, A. and TITS, J.: 'Groupes roouctifs', Publ. Math. IHES
27 (1965), 55-150.
AMS 1980 Subject Classification: 20D25, 20G99
BOREL SUMMATION MEI'HOD - A method for
summing series of functions, proposed by E. Borel [1].
Suppose one is given a series of numbers
(*)
let Sn be its partial sums and let S be a real number.
The series (*) is summable by the Borel method (B-
method) to the number S if
. -x 00 Xk _
x~ e k~O kf Sk - S.
There exists an integral summation method due to
Borel. This is the B' -method: If
a Uk
fe- u Lkk-
00 00
, du = S,
o k =0 .
then one says that the series (*) is summable by the
B '-method to the number S. For conditions under
which the two methods B and B' are equivalent, cf. [2].
The B-method originated in the context of analytic
extension of a function regular at a point. Let
n =0
be regular at the point 0 and let C be the set of all its
singular points. Draw the segment OP and the straight
line Lp normal to OP through any point PEe. The set
of points on the same side with 0 for each straight line
Lp is denoted by II; the boundary r of the domain II
is then called the Borel polygon of the function f (z),
while the domain II is called its interior domain. The
following theorem is valid: The series
is summable by the B' -method in II, but not in the
domain II which is the complement of II [2].
References
[I] BOREL, E.: 'Memoire sur les series divergentes', Ann. Sci. Ecole
Norm. Sup. (3) 16 (1899),9-131.
[2] HARDy, G.H.: Divergent series, Clarendon, 1949.
A.A. Zakharov
Editorial comments.
439
BOREL SUMMA nON METHOD

References sider the class of entire functions f (z) of exponential

[A1] BEEKMANN, W. and ZELLER, K.: Theorie der Limitierungsver- type ";;;0' for which
(ahren, Springer, 1970.
00

AMS 1980 Subject Classification: 40C05 j 1f (x) 12 dx < 00.

-00

BOREL SYSTEM OF SETS, B-system, generated by This class is identical with the class of functions f (z)
a system of sets M - The smallest (0', ~)-system of sets, that can be represented as
B (M), containing M. The sets belonging to B (M) are
called the Borel sets (cf. Borel set), or B-sets, generated fez) = .k-
v2'IT
!eiz1q,(t)dt,
-0

by the system M. For each ordinal number a<wI,

where WI is the initial ordinal number of cardinality ~I'
the Borel classes Ma are defined as follows: Mo=M;
M a consists of the unions if a is odd, and it consists of
the intersections if a is even, of sequences of sets
belonging to M fJ, f3<a. In such a case
B(M)= U {Ma: a<wd. The same Borel system of
sets B (M) is obtained if the union and intersection
above are interchanged. A Borel set belongs properly to
a class M a if it belongs to the class M a but does not
belong to M fJ if f3<a. (Sometimes one takes non-
intersecting classes, i.e. then the system
Ma \ U {MfJ: f3<a} is called a class.)
A.G. El'kin
AMS 1980 Subject Classification: 54A99
BOREL TRANSFORM - An integral transform of the
type 00
where cp(t)EL 2( -0',0').
References
[II BoREL, E.: Lefons sur les series divergentes, Gauthier-Villars,
1928.
[21 DZHRBASHYAN, M.M.: Integral transforms and representation of
functions in a complex domain, Moscow, 1966 (in Russian).
A.F. Leont'ev
Editorial comments. The statement at the end of the
article above is called the Paley - Wiener theorem.
References
[A1] BOAS, R.P.: Entire functions, Acad. Press, 1954.
AMS 1980 Subject Classification: 44A 10, 44A 15
BORSUK PROBLEM - One of the fundamental
problems in combinatorial geometry: Is it possible, for
any bounded set of diameter d>O in an n-dimensional
Euclidean space, to make a decomposition into not
more than n + 1 subsets with diameters smaller than a?

yet) = jf(z)e- ZI dz,

The problem was formulated by K. Borsuk [1] who
o noted that it was not possible to subdivide an n-
where f (z) is an entire function of exponential type. dimensional simplex and an n-dimensional sphere in Rn
The Borel transform is a special case of the Laplace into n parts of a smaller diameter. The problem has a
transfonn. The function y(t) is called the Borel positive solution for n = 2, 3, but only partial results
transform of f (z). If have been obtained for n > 3. Thus, for instance, the
problem has been positively solved for any bounded
fez) = n=O
~ an n
~-,z ,
n.
smooth convex set in Rn [2]. It has been proved that
the solution of Borsuk's problem can be reduced to the
then 00 case of sets of constant width. If a (F) is the smallest
yet) = ~ avt -(v + I);
v=o number of parts of a diameter smaller than d into
which a set FeRn can be subdivided, then the equality
the series converges for I t I >0', where 0' is the type of
a (F) = 3 is valid for a figure Fe R2 of diameter d if
f (z). Let D be the smallest closed convex set contain- and only if R2 contains a unique figure of constant
ing all the singularities of the function y(t); let
width d containing F [3]. This fact cannot be directly
K(q,) = mRe(ze- i </
generalized to the case n > 2. The Borsuk problem is
ZED
closely related to the illumination problem and to the
be the supporting function of D; and let h(cp) be the Hadwiger hypothesis, which is a generalization of the
growth indicator function of fez); then K(cp)=h( -cp). Borsuk problem in which Rn is replaced by a finite-
If in a Borel transform the integration takes place over dimensional normed space.
a ray argz =cp, the corresponding integral will converge
in the half-plane xcoscp+ysincp>K(-cp). Let C be a References
[IJ BORSUK, K.: 'Drei Siitze tiber die n -dimensionale euklidische
closed contour surrounding 15; then Sphare', Fund. Math. 20 (1933). 177190.
[2) GRDNBAUM, B.: Borsuk's problem and related questions, Proc.
fez) 1 . jy(t)e ZI dt.
= -2 Symp. Pure Math .. 7. Amer. Math. Soc .. 1963. pp. 271-284.
7Tle
[3) BOLTYANSKlI, V.G.: 'On decomposition of plane figures in
If additional conditions are imposed, other representa- parts of least diameter', Col/oq. Math. 21, no. 2 (1967), 253263
(in Russian).
tions may be deduced from this formula. Thus, con- P.S. Soltan

440
 BOSE - EINSTEIN STATISTICS

Editorial comments. and defines the probability of the particle distribution

References over the cells, characterized by the occupation numbers
[A1] BOLTYANSKlT, V.G. and GOKHBERG, I.Ts.: satze und Prob- NJ,N 2 ,
leme der Kombinatorische Geometrie, Deutsch. Verlag The most-probable distribution corresponding to a
Wissenschaft., 1972 (translated from the Russian).
given energy E and number of particles N:
AMS 1980 Subject Classification: 51 N20, 05899
E = ~(iNi' N = ~Ni' (2)

BOSE - EINSTEIN STATISTICS, Bose statistics -

is found from the extremum of (1) under the comple-
Quantum statistics used in systems of identical particles
mentary conditions (2). The corresponding average
with integral spin (0, 1, ... , in the unit
occupation numbers are
n=1.05XIO- 27 ergXsec). Presented by S. Bose and A.
Einstein in 1924. According to these statistics an arbi- -
ni
= KG = _1_ - I
e{J(,-I') ,
(3)
i
trary number of particles may be found in each quan-
tum state. It was shown by W. Pauli that the type of where p. is the chemical potential, 13= 1 / kT, k is the
quantum statistics is directly related to the particle Boltzmann constant (a universal constant
spin, since the population of particles with an integral k = 1.38 X 10- 16 erg / degree), and T is the absolute
spin obeys the Bose - Einstein statistics, while those temperature. The values of 13 and p. can be found from
with half-integral spin obey the Fenni - Dirac statistics. the conditions (2).
The state of a system of many particles is defined in The entropy of the system is defined as the logarithm
quantum mechanics by a wave function which, in the of the statistical weight (1) for the most-probable distri-
case of identical particles, may be symmetric with bution (3):
respect to the rearrangement of any pair of particles
(for particles with integral spin) or anti-symmetric (for
S = kIn W{n;} = k~Gi{niin [1+ ~i) +In(l+ni>}- (4)

particles with half-integral spin). For a system of parti-

Other thermodynamic functions may be derived from
cles obeying the Bose - Einstein statistics the states are
the entropy and the average energy.
described by symmetric wave functions. This is a
The usual approach to Bose - Einstein statistics con-
second, equivalent, formulation of the Bose- Einstein
sists in employing the large canonical Gibbs distribu-
statistics. Systems of a large number of particles obey-
tion for the probability Wn of the quantum level n of
ing Bose - Einstein statistics are known as Bose systems
the entire system being filled
(e.g. a Bose gas).
For an ideal quantum gas, i.e. for a system of identi-
cal particles with mass m which do not interact, in a
Wn = Q-I exp { - En;:}, (5)
cube of volume V = L 3 the quantum energy levels for
single particles are given by
where
Q = ~exp
{}
- En;: (6)
-L
2m'
(p -
is the sum over states, or partition function, and En are
where p are the eigen values of the momentum of the the energy levels of the entire system. Thus, for
individual particle: p = 2wnn / L, where n is a vector instance, for an ideal Bose gas
with integer (positive, negative or zero) components.
The quantum state of an ideal gas is defined by
En = ~..lC..np,
p2m
2
n = {np}, np=O, 1, ... ,

specifying the population of occupation numbers of the
levels {np }, where each np is the number of particles in
the one-particle state p. For Bose systems np =0, 1, ....
For large systems the energy levels are very densely
located and tend to form a continuous spectrum as
V ~ 00. Let the levels be grouped by small cells, con-
taining G; levels per cell. To each cell corresponds the
average energy ;, and the number G; is assumed to be
very large. The state of the system is defined by the
selection {N;}, where N; is the sum of np taken over the
levels of the cell. The statistical weight, that is, the
number of different particle distributions over the cells,
and from equations (5) and (6) formula (3) for the
occupation numbers, and formula (4) for the entropy
can be found, without the use of combinatorial
analysis. Such an approach is particularly important in
the case of non-ideal Bose systems, when the use of
Bose - Einstein statistics cannot be reduced to a simple
combinatorial problem. In such a case the requirements
of Bose- Einstein statistics can be met if for the Ham-
ilton operator H the representation of second quantiza-
tion is used, in which its action is defined on the space
of symmetric wave functions or on the space of occupa-
tion numbers. The sum over states will then be

H;r},
is
(I) Q = spexp {

441
BOSE-EINSTEIN STATISTICS
where N is the operator for the number of particles;
this sum may be used to find all the thermodynamic
functions of the Bose system.
References
[I] HUANG, K.: Statistical mechanics, New York, 1963.
[2] KUBO, R.: Statistical mechanics, North-Holland, 1967.
[3] LANDAU, L.D. and LIFSHITZ, E.M.: Statistical physics, Per-
gamon, 1980 (translated from the Russian).
[4] SCHRODINGER, E.: Statistical thermodynamiCS, Cambridge
Univ. Press, 1948.
[5] BOGOLYUBOY, N.N.: Lectures on quantum statistics, Gordon
and Breach, 1967 (translated from the Russian).
D.N. Zubarev
AMS 1980 Subject Classification: 81-XX, 82A 15
BOTI PERIODICITY THEOREM - A fundamental
theorem in K-theory which, in its simplest form, states
that for any (compact) space X there exists an isomor-
phism between the rings K(X)K(S2) and K(XXS2).
More generally, if L is a complex vector bundle over X
and P(LEB 1) is the projectivization of LEB 1, then the
ring K(P(LtB 1 is a K(X)-algebra with one generator
[H] and a unique relation ([H]-[I])([L][H]-[1])=O,
where [E] is the image of a vector bundle E in K(X)
and H -( is the Hopf fibration over P (L EB 1). This fact
is equivalent to the existence of a Thorn isomorphism
in K-theory for complex vector bundles. In particular,
P(ltBI)=XXS2. Bott's periodicity theorem was first
demonstrated by R. Bott [1] using Morse theory, and
was then re-formulated in terms of K-theory [6]; an
analogous theorem has also been demonstrated for real
fibre bundles.
Bott's periodicity theorem establishes the property of
the stable homotopy type of the unitary group U, con-
sisting in the fact that Q2 U,,-, U, where QX is the space
of loops on X, and "-' is weak homotopy equivalence,
in particular 7Tj( U) = 7Tj + 2 (U) for i = 0, I, ... , where 7Ti
is the i-th homotopy group. Similarly, for the orthogo-
nal group 0:
References
[I] BOTT, R.: 'The stable homotopy of the classical groups', Ann.
of Math. (2) 70, no. 2 (1959), 313-337.
[2] MILNOR,1.: Morse theory, Princeton Univ. Press, 1963.
[3] ATIYAH, M.F.: Ktheory: lectures, Benjamin, 1967.
[4] HUSEMOLLER, D.: Fibre bundles, McGraw-Hili, 1966.
[5] MOORE, J.e.: 'On the periodicity theorem for complex vector
bundles', in Sem. H. Cartan, 1959-1960.
[6] ATIYAH, M.F. and BOIT, R.: 'On the periodicity theorem for
complex vector bundles', Acta Math. 112 (1964),229-247.
A.F. Shchekut'ev
Editorial comments.
References
[A1) BOTT, R.: Lectures on K(X), Benjamin, 1969.
[A2) KAROUBI, M.: K-theory, Springer, 1978.
AMS 1980 Subject Classification: 55N15
BOUND VAlUABLE, bound occurrence of a variable -
442
A type of occurrence of a variable in a linguistic
expression. An exact definition for each formalized
language depends on the rules of formation of the
language. It is not possible to substitute objects in
place of a bound variable, since such a substitution
would lead to meaningless expressions. But the
replacement of a bound variable, where it occurs, by a
new (for a given expression) variable leads to an
expression with the same meaning. For example, in the
expressiOns
jf(x,y)dx, {x: f(x,y)=O},
x is a bound variable. Replacing x by a number leads
to a meaningless expression, whereas by writing z
everywhere instead of x one obtains an expression with
exactly the same meaning.
Bound variables always arise in applying to an
expression e with free occurrences of a variable x an
operator variable x (see Free variable). In the resulting
expression, all the occurrences of x in e that were previ-
ously free become bound. We mention below certain
operators that are often used (next to the operators
f ... dx and {x: . . . }), in which x is an operator
variable:
V'x( ... ), 3x( ... ), that is, the universal and
existential quantifiers;
f. ..... ... dx, that is, a definite integral with respect
to x;
~ " ' , that is, summation over x;
x
Ax( ... ), that is, a function of x the value of which
at x is .... Specific linguistic expressions can be sub-
stituted in place of the dots.
In real (non-formalized) mathematical texts it is pos-
sible to have a non-unique use for one and the same
expression; in this connection distinguishing a bound
variable in a given expression depends on the context
and meaning of the expression. In formalized languages
there is a formal procedure for distinguishing free and
bound occurrences of variables.
v.N. Grishin
AMS 1980 Subject Classification: 03-XX, 03005
BOUND VECfOR - See Vector.
AMS 1980 Subject Classification: 15A03
BOUNDARY - I) The boundary of a subspace A of a
given topological space X is the set of points such that
every neighbourhood of any point of it contains both
points from A and points from X\A. The accepted
notations include aA, Fr A, Frx A.
2) A synonym for the border of a manifold, such as
the border of a simplex.
A. V. Chernavskif
AMS 1980 Subject Classification: 54AXX, 57-XX
 BOUNDARY (IN THE THEORY OF UNIFORM ALGEBRAS)

BOUNDARY CONDmONS - Conditions imposed There are several known compactifications of a

on an (unknown) solution of a differential equation on simply-connected domain with the property of boun-
the boundary (or on a part of the boundary) of the dary correspondence under conformal mapping. His-
domain in which the solution is to be determined. torical precedence goes to the Caratheodory extension
Boundary conditions are usually formulated in terms of (see [1], and also [2]). It is the most intuitive and is
differential operators, though other types are also often used in the study of conformal and other map-
encountered. pings. The elements of the boundary thus obtained
A.P. Soldatov
were called prime ends by C. Caratheodory (see Limit
AMS 1980 Subject Classification: 34-01, 35-01 elements). A theory has been developed of boundary
correspondence under variable conformal mappings of
BOUNDARY COKRESPONDENCE, PRINCIPLE OF simply-connected domains (see [3]).
- A principle formulated in the following way. One says
References
that the principle of boundary correspondence holds [1] MYSHKlS, A.D. and SUVOROV, G.D.: 'Conformally-invariant
for a mapping f if ~e facts that f is a continuous mae: bicompact extensions of a plane simply-connected domain',
Soviet Math. Dokl. 14 (1973-1974), 1488-1491. (Dokl. Akad.
ping of the closure G of a domain G onto the clos~re D
Nauk SSSR 212 (1973), 822-824)
of a domain D and f is a homeomorphism of G \ G [2] CARATHEODORY, c.: 'Uber die Begrenzung einfach
onto 15 \ D imply that f is a topological mapping of G zusammenhangender Gebiete', Math. Ann. 73 (1913), 323-370.
[3] SUVOROV, G.D.: 'Prime ends of a sequence of plane regions
onto D. Thus, the principle of boundary correspon- converging to a nucleus', Amer. Math. Soc. Transl. Ser. 2 1
dence is in some sense converse to the boundary- (1955),67-93. (Mat. Sb. 33 (75) (1953), 73-100)
correspondence principle (cf. Boundary correspondence [4] MARKUSHEVICH, A.I.: Theory offunctions of a complex variable,
(under confonnal mapping)). 1-3, Chelsea, 1965-1977 (translated from the Russian).
[5] COLLINGWOOD, E.F. and LOHWATER, A.J.: The theory of clus-
If G and D are plane domains with Euclidean boun- ter sets, Cambridge Univ. Press, 1966.
daries homeomorphic to a circle and D is bounded, [6] SUVOROV, G.D.: Families ofplane topological mappings, Novo-
then the principle of boundary correspondence holds sibirsk, 1965 (in Russian).
for analytic functions f on G, i.e. f is a conformal map- [7] SUVOROV, G.D.: Metric theory ofprime ends and boundary pro-
perties of plane mappings with bounded Dirichlet integrals, Kiev,
ping of G onto D. In addition to that given above, vari- 1981 (in Russian).
ous other forms of the principle of boundary [8] IVANOV, O.V. and SUVOROV, G.D.: Complete lattices of
conformally-invariant compactifications of a domain, Kiev, 1982
correspondence are commonly used for conformal map-
(in Russian).
pings (see [1]). The principle of boundary correspon- B.P. Kufarev
dence has been verified for orientable mappings in Editorial comments. Standard English references on
Euclidean space (see [2]). boundary correspondence under conformal mapping and
prime ends are [A1] - [A3].
References
[I] LAVRENT'EV, M.A. and SHABAT, B.V.: Methoden der komplexen References
Funktionentheorie, Deutsch. Verlag Wissenschaft., 1967 [A1] AHLFORS, L.V.: Conformal invariants, McGraw-Hili, 1973.
(translated from the Russian). [A2] OHTSUXA, M.: Dirichlet problem, extremal length and prime
[2] KUDRYAVTSEV, L.D.: 'On differentiable mappings', Dokl. ends, van Nostrand Reinhold, 1970.
Akad. Nauk SSSR 95, no. 5 (1954), 921-923 (in Russian). [A3] POMMERENXE, c.: Univalent functions, Vandenhoeck and
[3] PINCHUX, S.I.: 'Holomorphic equivalence of certain classes of Ruprecht, 1975.
domains in C", Math. USSR-Sb. 111 (153), no. I (180), 67-94;
159. (Mat. Sb. 111 (153), no. 1 (1980), 67-94; 159) AMS 1980 Subject Classification: 30C35, 30040
B.P. Kufarev
BOUNDARY (IN THE THEORY OF UNIFORM
AMS 1980 Subject Classification: 30040, 54C05 ALGEBRAS) - A subset f of the space MA of maximal
ideals of a commutative Banach algebra A with an iden-
BOUNDARY CORRESPONDENCE (UNDER CON- tity over the field C of complex numbers, on which the
FORMAL MAPPING) - A property of a univalent con- moduli of the Gel'fand transforms a
of all elements
formal mapping f of a finitely-connected domain G a EA attain their maximum (cf. Gel'fand representa-
onto a domain D in the z-plane consisting of the fact tion). For example, one can set f=MA (the trivial
that f can be extended t? a hO!!leomorphism between boundary). Of interest are non-trivial boundaries with
certain compactifications G and D of G and D, respect- some sort of minimality property. There exists among
ively; that is, f induces a homeomorphism of the boun-
- -
daries G \ G and D \ D. For the ordinary (Euclidean)
the closed boundaries fCMA a minimal one aMA , that
is, a closed boundary such that aMA c f for every
boundaries aG and aD of G and D this property does closed boundary f; it is called the Shilov boundary. The
not always hold. For example, a conformal mapping of points of the Shilov boundary are characterized by the
a disc K induces a homeomorphism of aK and aD only property that for each neighbourhood V CMA of such
if aD is homeomorphic to a circle. a point ~ and every >0, there exists an element a EA

443
BOUNDARY (IN THE THEORY OF UNIFORM ALGEBRAS)
a a
for which max I I = I and I I < outside V. The
points ~EoMA constitute 'the most stable' part of the
set MA of maximal ideals: If B is a commutative
Banach algebra containing A as a sub algebra, then the
maximal ideals (multiplicative functionals) correspond-
ing to such points can be extended to maximal ideals
(multiplicative functionals) of B, whereas for maximal
ideals not belonging to the boundary oMA such an
extension is, in general, not possible. This situation is
analogous to the stability of the boundary of the spec-
trum of a bounded linear operator on a Banach space.
A typical example is the algebra A consisting of the
functions continuous on the disc I A I .;;;; I that are ana-
lytic for I A I < 1. In this case MA can be identified
with the closed disc and oMA with its topological boun-
dary; maximal ideals corresponding to interior points
of the disc cannot be extended to maximal ideals of the
algebra of all functions continuous on the boundary in
which A is naturally included (according to the max-
imum principle), while the maximal ideals correspond-
ing to the boundary points can be extended.
As in the case of algebras of analytic functions there
is a local maximum-modulus principle for general com-
mutative Banach algebras: If V is an open subset of the
space M A , then
max{ I a(~) I : ~E V} = max{ I a(0 I : ~E(aMA n V) U aV}
for all a EA, where V is the closure and dV is the topo-
logical boundary of V in M A Roughly speaking, this
means that a local maximum point is necessarily a glo-
bal maximum point (possibly of another element).
The notion of a boundary is used in the study of uni-
form algebras, that is, closed sub algebras A of the alge-
bra C(X) of all continuous functions on a compactum
X that separate points and contain the constants. In
this situation OMA CXCMA. There exists for each
point ~EMA a representing measure, concentrated on
the Shilov boundary, that is, a probability measure ~
such that
for all a EA (this holds for arbitrary commutative
Banach algebras). In the simplest case (described
above) of the disc and the algebra of analytic functions,
the latter formula reduces to the Poisson formula. The
representing measure J1. is in general not unique. For
points belonging to the same Gleason part (see Algebra
of functions), the measures J1. can be chosen absolutely
continuous with respect to one another, which, under
certain extra conditions of uniqueness type for the
representing measures, makes it possible to equip the
Gleason parts of the space of maximal ideals with a
one-dimensional analytic structure compatible with the
algebra. Every point of the Shilov boundary of a uni-
444
form algebra constitutes a one-point Gleason part;
however, the converse is in general false.
The equality oMA =X =MA for uniform algebras is a
simple necessary condition for the coincidence of A
with C(X). Without extra hypotheses this is not a suffi-
cient condition for the above coincidence even in the
case of algebras A =R(X) of uniform limits of rational
functions on a compactum X in the plane.
Let X be a metrizable compactum and let A be a
uniform algebra on X. Then there exists a minimal
boundary among all the boundaries: ooMA. The closure
of the minimal boundary is the same as oMA However,
OOMA is, in general, not closed; an example is the
subalgebra of all analytic functions inside the disc
I A I :s;;; I for which f(O) =f(1). The boundary OOMA is
the same as the set of peak points with respect to A,
that is, points ~EMA for which there exists an a EA
with I am I > I a(1/) I for all1/*~. On the other hand,
a condition for a point ~ to belong to doMA that is for-
mally considerably weaker is known. Namely, if the
point ~EMA is such that for some O<c< I and d~ I
there exists for each neighbourhood V of ~ an element
a in A for which a(~)= 1, max I a I =d and I a(1/) I.;;;;c
for 1/ft V, then ~EooMA' The abstract Poisson formula
has been strengthened in the following way: There
exists for any point ~EMA a representing measure con-
centrated on OOMA (doMA being, moreover, a Gil-set).
In this form it has been successfully applied in certain
problems of approximation theory. The points
~EooMA are characterized by the property that for
these points, the measure is unique and is the same as
the Dirac ~-measure, that is, the minimal boundary is a
special case of the Choquet boundary.
For the algebra A =R(X) of uniform limits of
rational functions on a compactum X in the plane, A
and C(X) coincide if and only if X=MA and OOMA
coincide. This is not true for arbitrary uniform alge-
bras: There exists a uniform algebra A distinct from
C(MA) for which MA is metrizable and for which each
point ~EMA is a peak point (that is, OOMA =MA)'
There also exists a uniform algebra for which all the
Gleason parts are trivial (singletons), but for which
even the Shilov boundary is a proper part of the space
of maximal ideals.
One of the algebraic generalizations of the notion of
the Shilov boundary is the following. Let olMA =MA
and let I be the closed ideal generated by a set of ele-
ments al,'" ,an~1 EA, n;;;.2. The space MAil can be
naturally identified with the set of common zeros of the
a
functions I, . . . ,an ~ I . The closure of the union of
the boundaries OMA I lover all ideals I generated by
n - 1 elements is denoted by dnMA . For example, for
the algebra of continuous functions in the polydisc
I Al I:s;;; 1, I A2 I :s; ; 1 that are analytic in its interior,
 BOUNDARY-LA YER THEORY

a]MA is the skeleton I A] I = I A2 1=1, a2M A is the
topological boundary I A] I = 1, I A2 I ~ 1; I A] I ~ 1,
I A2 I = 1, and a3MA = MA These generalizations are
useful when proving theorems on multi-dimensional
analytic structures in spaces of maximal ideals.
Notions of boundaries (or functional boundaries)
similar to the ones described above are encountered in
the theory of analytic functions (the Bergman boun-
dary), probability theory (the Martin boundary) and in
a number of other branches of mathematics. In this
connection the initial set is not necessarily assumed to
be an algebra or a ring.
References
[1) BASENER, R.: 'A generalized Shiloy boundary and analytic
structure', Proc. A mer. Math. Soc. 47, no. I (1975), 98-104. problems) in subregions where there is a substantial
[2) BROWDER, A.: Introduction to function algebras, Benjamin,
1969.
[3) GAMELIN, T.W.: Uniform algebras, Prentice-Hall, 1969.
[4) GELFAND, I.M. [I.M. GEL'FAND), RAIKOV, D.A. [D.A. RAiKOV)
and SCHlLOW, G.E. [G.E. SHILOV): Kommutative Normierte
Ringe, Deutsch. Verlag Wissenschaft., 1964 (translated from
the Russian).
[5) GONCHAR, A.A.: 'The minimal boundary of the algebra A (E)"
Izv. Akad. Nauk SSSR Ser. Mat. 27, no. 4 (1963),949-955 (in
Russian).
[6) PHELPS, R.R.: Lectures on Choquet's theorem, v. Nostrand,
1966. E.A, Gorin represented as the sum of two expansions. The outer
Editorial comments. Although the concept of a boun-
dary layer has its origin in fluid dynamics, the recent
mathematical development of this concept, as well as its
applications, are in the field of singular perturbations. For a
partial review one may consult [A1].
References
[A 1] EcKHAus, W.: 'Boundary layers in linear elliptic singular per-
turbation problems', SIAM Review 14 (1972), 225-270.
AMS 1980 Subject Classification: 76D10
BOUNDARY-LAYER. THEORY - An asymptotic
approximation of the solution of boundary value prob-
lems for differential equations containing a small
parameter in front of the highest derivative (singular
effect from the terms containing the highest derivatives
on the solution. The boundary-layer phenomenon
arises in narrow zones near the parts of the boundary
on which there is a difference in the numbers of boun-
dary conditions for the initial problem and the degen-
erate one (with the small parameter taking the value
zero), as well as near the surfaces of discontinuity of
the solution of the degenerate problem.
The solution of the singular problem may be

Editorial comments. Much of this article, in particular expansion is determined by the method of the small
the definition of boundary, is relevant and valid for general parameter with the part Bo of the boundary conditions
commutative Banach algebras. B = ~ 0 +B ]. The inner expansion decreases rapidly
AMS 1980 Subject Classification: 46J20 outSIde the boundary layer and is usually sought for as
polynomials in powers of (. To determine these, the dif-
ferential equations are transformed by variables that
BOUNDARY LAYER - A region of large values of
depend on ( and that stretch the subregions of the
the gradient of a function, in particular, in hydro-
boundary layer. The equations of the boundary layer
dynamics it is a region of the flow of a viscous liquid
(gas) the transversal thickness of which is small in com- are derived by equating to zero the coefficients at vari-
parsion with its longitudinal dimensions and which is ous powers of ( after substituting the polynomials into
the transformed equations. To these one adds condi-
produced at the surface of a solid or at the boundary
tions B]. The error in the outer expansion, after it has
between two liquid flows with different velocities, tem-
peratures or chemical compositions. A boundary layer been found, indicates the necessary change of variables.
is characterized by a sharp change in the velocity in the In solving complicated applied problems, the necessary
transversal direction (a shear layer) or a sharp change c~ange ~f variables can be clarified on the basis of phy-
SICal esttmates for the terms in the initial equations and
in the temperature (a thermal, or temperature, boundary
the simplifications corresponding to them. This change
layer), or else in the concentrations of the individual
should eliminate the highest derivatives in front of f.
chemical components (a diffusion, or concentration,
To solve the problem it is necessary to determine where
boundary layer). The concept of a boundary layer and
the boundary layers lie and how the conditions B are to
the term itself were introduced by L. Prandtl (1904) in
connection with the solution of a boundary value prob- be separated into Bo and B]. A notable feature of this
lem for non-linear partial differential equations in the solution is that hyperbolic equations for the outer
hydrodynamics of viscous liquids. The needs of aviation expansions and parabolic equations for the inner
have led to the development of a boundary-layer theory e~pansions may correspond to initially elliptic equa-
in aerohydrodynamics. In the mid-twentieth century, nons.
the mathematical boundary-layer theory developed, as In the theory of systems of ordinary differential
well as applications in the theory of heat transfer, diffu- equations
sion, processes in semi-conductors, etc. XI = f (x, y, t), YI = g(x, y, t), O';;;;;t';;;;; T, (1)
Yu.D. Shmyglevskif where x, f are m-dimensional and y, g are n-

445
BOUNDARY-LAYER THEORY
dimensional vector functions, an existence and unique-
ness theorem for the solution of the Cauchy problem
under the conditions x (0, ()=x o and yeO, ()=yO and
under certain properties of j and g has been proved.
One has also derived properties of the solution as (~O
(see [1]). In the case of a boundary value problem for
(1) under the conditions
x(O,) = xo, YI(O,) = yy, Y2(T, ) = y~,
where the sum of the numbers of components of the
vectors Y I and Y2 is n, there exist, in general, boundary
layers in neighbourhoods of the ends of the segment
[0, T]. An algorithm has been constructed for finding
asymptotics of the solution for this problem. Under
certain properties of j and g it has been shown that a
solution exists and is unique, and it has been estimated
(see [3]). If a solution to the limit equation g=O is not
unique with respect to y, one can construct an internal
boundary layer (in a neighbourhood of 7',0<7'<1) that
separates regions with different solutions of the limit
equation. An algorithm has been constructed for a par-
ticular type of integro-differential equation giving an
asymptotic expansion in ( for the problem with initial
conditions, and some features of the behaviour of the
solution have been examined.
In the case of linear ordinary differential equations
(L+(M)x=j(t), O~t~l, where Land M are differen-
tial operators, with boundary conditions B + B I, one
can distinguish the class of problems the solutions of
which contain boundary layers, and the concept of reg-
ular degeneracy (the solution to the limit equation
enables one to satisfy the conditions B 0, while the
asymptotic solution for the boundary layer enables one
to satisfy B 1) can been introduced (see [4]). An iterative
process has been constructed for the asymptotic
representation of the solution, and estimates have been
given for the residual terms in the expansions.
It has been shown (see [5]) in the theory of boundary
layers of a general non-linear second-order ordinary
differential equation, subject to certain assumptions,
that the solution of the first boundary value problem is
made up of an external solution, a boundary layer and
a residual term that, together with its first-order deriva-
tive, is of order t: on the segment.
Studies have been made on the behaviour of solu-
tions of boundary value problems of basic types for a
linear partial differential equation of the form
ilu + A (x, y)ux + H(x, Y)Uy + C(x, y)u = f(x, y),
where D. is the Laplace operator, in a region D with
boundary S. Conditions on the functions A, B, C, and
f, the boundary S, and the functions a and cp of the
points P on S appearing in the boundary condition
Un +a(P)=</i,P) have been given such that u in D US
tends uniformly to the solution of the limit equation
446
with this boundary condition on a certain part of S
(absence of a boundary layer) [6].
For an elliptic second-order equation in a region D
with boundary S, using the example of two indepen-
dent variables
t:(a (x, y)uxx +2b (x, y)uxy +c(x, y)Uyy +
+d(x,y)ux +e(x,y)Uy +g(x, y)u]+ux -
-h(x,y)u = I(x,y), h~a2>0,
iterative processes have been constructed solving the
problem with the condition u =0 on S, theorems have
been proved on the structure of the expansion of u with
respect to t: and estimates have been made of the resi-
dual term in this expansion [4]. Similar results have
been obtained for equations of higher orders.
A method has been devised [7] for combining asymp-
totic expansions for the equation
ilu -a(x, y)Uy = I(x,y)
in a rectangle with u given on the boundary.
Research on boundary-layer theory for non-linear
partial differential equations is related mainly to
aerohydrodynamics and is based on the N avier - Stokes
equations or generalizations of them. Practical require-
ments have led to the development of the mathematical
theory and to methods of handling various problems.
Below only laminar flows are considered (see [8] - [10]).
The Navier-Stokes equations in the case of the
hydrodynamics of planar (k =0) and axi-symmetric
(k = 1) flows of an incompressible fluid with constant
viscosity coefficient pare:
(Ir"U)~+('l/kV~ = 0, UT+~U~+V~~ = ~ilU-P~'} (2)
VT+UV~+VV~ - ~ilv P~,
where t:=Re- 1j2 , Re=wX /p is the Reynolds number,
which is represented in terms of the characteristic
values of the velocity w and the linear dimension X.
The solutions are defined by the boundary conditions
on the closed boundary S of a region D, where on the
solid contour r one has the conditions U=O,
v=vo(~, H(~ for T/=H(~), where u and v are the
tangential and normal components to r of the vector
(u, v). The initial values of u, v and p are given on
DUS.
For t: smail, in a first approximation the asymptotic
solution is composed of the solution of (2) for =0
with some of the conditions on S (only the condition
v=vo is imposed on f) and the solution of the
boundary-layer equations. The equations for the
dynamic boundary layer are derived on the assumption
that the conditional thickness 8 of the boundary layer
and the value of v have orders 8"" X t:, V ""Wt:, and that
the terms on the left-hand sides of the latter equations
in (2) are of the order of the terms containing ~. Intro-

duction of the variables t =-r, x =g,y =1/ / f. and
V=v / f. leads, as f.~O, to the Prandtl equations (cf.
Prandtl equation):
(rku)x+(rkV)y = 0, Ut+uux+V~ = Uyy-Px.
Py = 0, O";;;;t,,;;;;T, O";;;;x,,;;;;Xo, O";;;;y<oo,
with the conditions
u\t=o = uo(x,y), v\t=o = vo(x,y), u\y=o = 0,
V\y=o = vo(t, x),
u _ W(t, x), px = - WWx - w" asy_oo, u \x=o = 0,
where r is the distance from the symmetry axis for
k = 1 and W (x) is a known function. These equations
and conditions apply for any curvilinear contour with
radius of curvature much larger than ~. In the latter
case, x and y are the coordinates along the contour and
along the normal to it.
If W is constant, the problem reduces to a boundary
value problem for an ordinary differential equation.
There are also other classes of analogous solutions.
Conditions under which solutions of boundary-layer
problems exist are known; one has investigated the
problem of uniqueness and stability of the solutions as
well as how they result from the solutions for station-
ary cases [11]. Solutions have been constructed by the
method of straight lines, and it has been shown that
they converge.
Boundary-layer equations for a compressible liquid
can be derived from the equations for the flow of a
viscous and heat-conducting gas; they are much more
complicated than (3). Their number is also larger.
There is an integral transformation that simplifies these
equations in the general case and reduces them to (3)
when the Prandtl number Pr=cp / K= 1, where cp is
the heat capacity of the gas at constant pressure and K
is the coefficient of heat conductivity [12]. Several
modifications of the transformation exist. In the general
case, the boundary-layer equations describe so-called
natural convective flows. If p is independent of the tem-
perature and if the Archimedean force is negligible,
then the energy equation splits off from the system of
boundary layer equations and one speaks of forced
convective flow. The energy equation determines the
thermal boundary layer, whose thickness differs from ~.
Boundary layers also arise in zones separating flows
with different characteristic velocities. Shock waves are
also boundary layers.
A distinct class of two-dimensional boundary-layer
problems is associated with flows in rotating axi-
symmetric plates and bodies.
Not only have methods been developed for solving
non-stationary problems, but also problems in which W
is periodic, when there is stepwise motion from a state
of rest, problems with accelerated motion, for the boun-
BOUNDARY-LAYER THEORY
dary layer behind a shock wave, and problems with a
variable temperature over the surface of the solid
around which the flow takes place have been solved.
(3) In the boundary-layer theory for three-dimensional
flows, methods for obtaining a solution have been
developed and cases in which the equations simplify
have been studied. The boundary-layer equations for a
sliding cylindrical wing of infinite span are analogous
to the equations for a two-dimensional boundary layer.
Approximate solutions have been obtained for the
problem of a boundary layer on a rotating cylindrical
propeller blade and on a rotating cylinder in a skew
flow, as well as for the boundary-layer problem near
the line of intersection between two planes.
These researches on boundary layers in aerohydro-
dynamics relate to a first approximation in boundary-
layer theory. Higher approximations enable one to
examine the interactions of boundary layers with the
external flow, and to make calculations for moderate
values of R [13].
The stability of boundary layers enables one to deter-
mine the limits of applicability of the theory. There are
studies [14] based on methods of small perturbations
with periodic and local initial perturbations. In the case
of two-dimensional flows, the analysis of three-
dimensional perturbations reduces in linear approxima-
tion, on the basis of Squire's theorem, to a two-
dimensional analysis with an altered value of P. Non-
linear analysis applied to stability loss shows that longi-
tudinal vortices occur.
Physical generalizations of problems in the theory of
boundary layers (see [15] - [17]) are related to research
on multi-phase flows, to the use of real equations of
state and transfer coefficients (a complication of the
equations), to the study of non-equilibrium flows with
diffusion (an extension of the system of equations,
which then becomes of parabolic-hyperbolic type), to
taking into account ablation of the surface around
which the flow takes place (a complication of the boun-
dary conditions and it becomes necessary to consider
the thermal conduction in the solid), and also to con-
sidering radiation transport (integro-differential equa-
tions).
A further development of boundary-layer theory in
aerohydrodynamics has been obtained in the study of
flows that do not satisfy Prandtl's assumptions, where
solutions are obtained as multi-layer asymptotic expan-
sions. The causes of the complexity in the structure of
the solution are that there are additional small parame-
ters in the boundary conditions (for example, because
the radius of curvature of the contour around which the
flow takes place is small), there may be singular points,
lines or surfaces in the first approximation of the
theory, and bifurcation of the solution is possible.
447
BOUNDARY-LAYER THEORY
1bis class also contains flows around points of
separation or connections of the boundary layers to the
contour around which the flow takes place and around
the points of incidence of shock waves on the boundary
layers. The characteristic solutions have ([18] - [20]) a
three-layer structure. The external solution defines the
potential flow perturbed by the boundary layer and is
described by equations in perturbations. The middle
layer of order X{ can be described by the equations for
non-viscous vortex flows with Py =0, and it receives gas
from the main part of the preceding boundary layer.
The equations for the third (thinnest) layer nearest to
the wall can be derived on the assumption that its
length and thickness are, correspondingly, of order
X~/4 and X~/4. The following variables are intro-
duced in the stationary case:
x= ~-3/4, Y = 'I}(- 5 / 4 ,
U = u(-I/4, V = v(-3/4, P = p(-1/2.
1bis leads, as {~O, again to (3) in which u is replaced
by U and p by P. 1bis structure for the solution of a
problem is characteristic for a wide class of flows with
small perturbations.
Many studies have been made on the dynamics of
flows for small {, where the pressure in the outer super-
sonic flow varies considerably over short distances
(M> 1, M = w / c, where w is the velocity of the gas
and c is the velocity of sound). This includes problems
on calculating flows around contours of large local cur-
vature and flow attachment to the surface of the body.
In these cases Py=l=O in the middle layer of the three-
layer scheme for the solution.
The class of problems in which the perturbed solu-
tion occupies a finite region has also been studied.
These solutions occur when there is moderate or strong
interaction of the boundary layer with an external
hypersonic flow (M~oo) or in supersonic flow around
a body of finite length in which there are intense gas
injections through the surface (vo >0). In these cases,
the pressure on the outer boundary of the boundary
layer is determined from the solution to the complete
problem. The perturbation from the trailing edge of the
body propagates upstream, which is caused by the
non-uniqueness of the solution in a neighbourhood of
the leading edge.
References
(1) TnrnONOV, A.N.: 'Systems of differential equations containing
small parameters in front of the derivatives', Mat. Sb. 31, no. 3
(1952), 575-586 (in Russian).
(2) W ASOV, W.: Asymptotic expansions for ordinary differential
equations, Interscience, 1965.
[3] VASIL'EVA, A.B. and Buruzov, V.F.: AsymptotiC expansions of
solutions to singularly perturbed equations, Moscow, 1973 (in
Russian).
[4A] VISHIK, M.I. and LYUSTERNIK, L.A.: 'Regular degeneracy and
boundary layers for linear differential equations with a small
448
parameter', Uspekhi Mat. Hauk 12, no. 5 (1957), 3-122 (in Rus-
sian).
[4B] VISHIK, M.1. and LYUSTEllNIK, L.A.: 'Solutions to some prob-
lems on perturbations in the case of matrices and non-self-
adjoint differential equations', Uspekhi Mat. Hauk 15, no. 3
(1960),3-80 (in Russian).
[5] CoLE, J.: Perturbation methods in applied mathematics, Blais-
dell, 1968.
[6] OLEiNIK, O.A.: 'on equations of elliptic type with a small
parameter in front of the highest derivatives', Mat. Sh. 31, no.
1 (1952), 104-117 (in Russian).
[7] !L'IN, A.M. and LEuKovA, E.F.: 'The method of matched
asymptotic expansions for the equation
/;u-a(x,y)u,=f(x,y) in a rectangle', Mat. Sh. 96, no. 4
(1975), 568-583 (in Russian).
[8] KocHIN, N.E., KmEL', LA. and ROZE, N.V.: Theoretical
hydromechanics, 2, Moscow, 1%3 (in Russian).
[9] LoiTSYANSKIi, L.G.: Laminar boundary layers, Moscow, 1962
(in Russian).
(10) ScHLICHTING, H. VON: Grenzschicht-Theorie, Karlsruhe, 1951.
[II] OLEINiK, O.A.: 'Mathematical problems of boundary-layer
theory', Uspekhi Mat. Haule 23, no. 3 (1968), 3-65 (in Russian).
[12] DoRODNITSYN, A.A.: 'Boundary layers in an incompressible
gas', Prikl. Mat. i Mekh. 6, no. 6 (1942), 449-486 (in Russian).
[13] DYKE, M. VAN: Perturbation methods inftuid mechanics, Para-
bolic Press, 1975.
[14] BETcHov, R. and CluMiNALE, V.O.: Stability ofparallelftows,
Acad. Press, 1967.
[15] Soo, S.: Fluid dynamics ofmultiphase systems, Blaisdell, 1971.
[16] DoRRANCE, W.H.: Viscous hypersonicftow, McGraw-Hill, 1962.
[17] KAyS, W.M.: Convective heat and mass transfer, McGraw-Hill,
1971.
[18] NEiLAND, V.YA.: 'On the theory of separation of laminar
boundary layers in supersonic flow', Izv. Akad. Haule SSSR
Mekh. Zhidk. i Gaza, no. 4 (1969), 53-57 (in Russian).
[19] NEiLAND, V.YA.: 'Asymptotic problems in the theory of
viscous flows', Trudy TsAGI, no. 1529 (1974) (in Russian).
[20] STEWARTSON, K.: 'Multistructured boundary layers on flat
plates and related bodies', Adv. App. Mech. 14 (1974), 145-239.
Yu.D. Shmyglevskii
Editorial comments. The article above presents a one-
sided view of the subject and is in many aspects not up-to-
date. Furthermore, very substantial contributions in the
Western literature are neglected. Mathematically boundary
layers arise in problems of singular perturbations and are
treated by methods of asymptotic analysis. In linear and
non-linear, ordinary and partial, differential equations much
more has been accomplished than described above in the
first part. Overviews that are partially complementary can be
found in [A1] - [A3].
From the point of view of applications boundary-layer
theory is described above as a branch of fluid- and gas-
dynamics, where it originated. However, boundary layers
and transition layers also occur in many other fields, such
as for example combustion, geophysics, free boundary
problems, or epidemology. Samples and further literature
can be found in [M], [A5).
Systems of (ordinary) differential equations of the type (1)
above are often referred to as singular perturbations (cf.
also the editorial comments to Perturbation theory).
For an idea of the role of singular perturbations, boundary
layers and multiple time scales in control and system theory,
cf. [A6).
References
[A 1] LIONS, J.L.: Perturbation singu/ieres dans les problemes aux

limits et en contrOle optimal, Lecture Notes in Math., 323,
Springer, 1973.
[A2] EcmAus, W.: Asymptotic analysis of singular perturbations,
North-Holland, 1979.
[A3] CHANG, K. W. and HOWES, F.A.: Nonlinear singular perturba-
tion phenomena: theory and application, Springer, 1984.
[A4] BRAUNER, C.M., GAY, B. and MATHIEu, J. (EDS.): Singular
perturbations and boundary layer theory, Lecture Notes in
Math., 594, Springer, 1977.
[A5] VERHULST, F. (ED.): Asymptotic analysis II, Lecture Notes in
Math., 985, Springer, 1983.
[A6] KOIWTOVIC, P. (ED.): Singular perturbations and time scales
in modelling and control of dynamical systems, Univ. of Illi-
nois at Urbana-Champaign, 1980.
AMS 1980 Subject Classification: 34E15, 35-01,
76010
BOUNDARY (OF A MANIFOLD) - The subset of
the closure M n of an (open) n-dimensional real mani-
fold M n for which a neighbourhood of each point is
homeomorphic to some domain wn in the closed half-
space of Rn , the domain being open in R'!t- (but not in
Rn). A point a EM n corresponding to a boundary point
of wn CR'!t-, i.e. to an intersection point of wn with
the boundary of R'!t-, is called a boundary point of Mn.
A manifold having boundary points is known as a man-
ifold with boundary. A compact manifold without boun-
dary is known as a closed manifold. The set of all
boundary points of M n is an (n -I)-dimensional mani-
fold without boundary.
M.I. Voitsekhovskii
Editorial comments.
References
[A1] HIRSCH, M.W.: Differential topology, Springer, 1976.
AMS 1980 Subject Classification: 57-XX
BOUNDARY PROPERTIES OF ANALYTIC FUNC-
TIONS - Properties of analytic functions that are
displayed as the function approaches the boundary of
its domain of definition.
I t can be said that the study of boundary properties
of analytic functions, understood in the widest sense of
the word, began with the Sokhotskii theorem and the
Picard theorem about the behaviour of analytic func-
tions in a neighbourhood of isolated essential singular
points (cf. Essential singular point), which were
obtained in the second half of the nineteenth century.
The terms relevant to this approach to the study of
boundary properties of analytic functions - which is
now called the theory of prime ends and cluster sets
(cf. Limit elements) - first appeared in a course given
by P. Painleve in 1895. The dissertation of P. Fatou
(1906) is the first systematic study of certain boundary
properties of analytic functions in a neighbourhood of
the continuous boundary of their domain of definition.
The theory of boundary properties made considerable
BOUNDARY PROPERTIES OF ANALYTIC FUNCTIONS
advances in the first third of the twentieth century,
owing to the work of several scientists; it resumed its
rapid advance in the second half of that century,
accompanied by the appearance of new ideas and
methods, novel directions and objects of study. Its
development is closely connected with various fields of
mathematical analysis and mathematics in general, first
and foremost with probability theory, the theory of har-
monic functions, the theory of conformal mapping,
boundary value problems of analytic function theory,
potential theory, value-distribution theory, Riemann
surfaces, subhartnonic functions and function algebras.
The theory of boundary properties of analytic functions
is closely connected with various fields of application of
mathematics by way of boundary value problems.
Since the study of boundary properties is connected,
in the first place, with the geometry of the boundary f
of the domain of definition D of an analytic function
f (z) in one complex variable z, three main approaches
can be distinguished in the theory of boundary proper-
ties of analytic functions.
a) The study of the behaviour of f (z) in a neigh-
bourhood of an isolated boundary point aEf. The
most important case is that of an essential singular
point a, which is dealt with in the theorems of
SokhotskiI, Picard, Julia, and Iversen (cf. Sokhotskii
theorem; Picard theorem; Julia theorem; Iversen
theorem).
b) The study of the behaviour of f (z) in the case
when the boundary f is an everywhere-discontinuous
set. V.V. Golubev's dissertation Single-valued analytic
functions with a perfect set of singular points (1916, cf.
[1]) was of great importance in this connection.
c) The study of the behaviour of f(z) when the
domain D is bounded by a continuous closed curve f
and, in particular, by the unit circle.
Cases a) and c) are, in a sense, extreme cases, while
case b) is intermediate. Case c), which is discussed
below, has been the subject of most intense study.
Let an analytic function f (z) be defined in a finite
simply-connected domain D of the complex z-plane
bounded by a rectifiable Jordan curve f. The following
problems are fundamental in the classical approach to
the study of boundary properties of analytic functions.
1) The problem of the existence of boundary values, i.e.
the question under which conditions and in which sense
the boundary values of f (z) exist as the point z
approaches f. This problem, as well as the succeeding
ones, can be formulated in a different manner as the
problem of identifying sufficiently extensive classes of
analytic functions in D that have, in some sense, boun-
dary values for sufficiently large sets of points on f.
2) The problem of boundary representation of f (z), i.e.
the problem under which conditions and with what
449
BOUNDARY PROPERTIES OF ANALYTIC FUNCTIONS
kind of analytic apparatus the dependence of f(z) on
its boundary values on r can be expressed. Oearly, this
apparatus will be different for different classes of ana-
lytic functions.
3) The uniqueness problem, or the problem of the pro-
perties that a set E C r should have such that two ana-
lytic functions of a given class coincide in D if their
boundary values on E are identical.
The first result in the solution of the existence prob-
lem is the theorem of Fatou (1906): If an analytic func-
tion is bounded in the unit disc D={z: I z I <I},
I f(z) I :s;;;.M, then radial boundary, or limit, values
f(ei8)=limHl_of(reiB) exist almost everywhere with
respect to the Lebesgue measure on the unit circle
r ={z: I z I = I}. It can be shown that, under these
conditions, not only radial, but also angular boundary
values, or boundary values along all non-tangential
paths, exist almost everywhere on r. This means that,
for almost-all points e iB Er, f(z) tends to a definite
limit f(e iB ) as z tends to the point e iB while remaining
-+
within an arbitrary fixed angle
~(eiB,() = {I z I<I} n {I arg(eiB-z) 1<;
>0, of width 7T-2<7T, with apex at the point e iB ,
bisected by the radius drawn through the point e iB .
Fatou's theorem cannot, in a certain sense, be
improved upon; it was in fact shown by N.N. Luzin
(1919) that for any set E Cr of measure zero on f
there exists a bounded analytic function f (z) not hav-
ing radial limits on E.
The class of bounded analytic functions in a domain
D is denoted by B(D) or HOO(D). Following the results
of Fatou, the next problem appeared to be the generali-
zation of his theorems to wider classes of functions.
One distinguishes between the following basic classes of
functions in the unit disc D, which are related by
proper inclusions:
A(D) C B(D) = Hoo(D) C HP(D) C N*(D) C N(D).(I)
The class A (D) is the class of single-valued analytic
functions in D that are continuous in the closed
domain D U r=D.
The classes HP(D), for all positive numbers p, are
defined by the condition
I 2" }I /P = C(f,p)
II I lip = O<r<:1
sup { -2 jl/(rei<P) IPdq,
'IT 0
< (2)
< +00.
The proper inclusions H oo CHP, CHP, are valid for
any O<p I <P2 < + 00. The classes HP were first intro-
duced by G.H. Hardy (1915), and are often named
Hardy classes. If I ~ < 00, one can introduce the norm
(2) on HP, and the norm
450
II I II = II I II B = ~~E II (z) I
00
on Hoo; the classes HP, I:s;;;.p:s;;;. + 00, have
a natural
structure of a vector space, and become Banach Hardy
spaces. If 0<p<1, it is only possible to introduce the
metric ppCf, g)= II f - g II~ on HP, which converts the
latter into a complete metric non-normable space. The
class of bounded analytic functions B =H OO is con-
tained in any class HP, P >0.
The class N(D) of certain meromorphic functions
f (z) in the unit disc D is said to be the class of func-
tions of bounded characteristic; it was introduced in
1924 by R. Nevanlinna. The class N(D) can be
characterized as the set of meromorphic functions f (z)
in D that can be represented as the ratio of two
bounded regular functions fl(z) and h(z) in D,
f(z)=fl(z)/ h(z).
All regular functions f(Z)EN(D) form a subclass
N*(D), andf(z)EN*(D) if and only if the condition
1 2"
sup -2 fln+ I/(rei<P) I dq, = C(f) < + 00 ' (3)
O<r<:1 'IT 0
where In+ x=Inx if Inx;;;;o.O and In+ x=O if Inx<O is
fulfilled. The class N* (D) contains all classes HP,
O<p:s;;;.+oo.
The classes HP have the following generalization. Let
#,t) be a strictly-convex function for - oo<t< + 00, i.e.
a non-negative, convex, non-decreasing function such
that #,t)/t~+oo as t~+oo. The class H",(D) is
then defined by the condition
I 2"
sup 1/1-1 { -2 jl/;{ln Il(re'<P) I)dq,
O<r<:1 'IT 0
. }
= C(f,I/I) < (2')
< +00,
compare with condition (2), where 1/J(t)=eP1 .
The principal result obtained for the problem of the
existence of boundary values in the case of the unit disc
D states that each meromorphic function f(z) of
bounded characteristic in D has angular boundary
values f(e iB ) almost everywhere on f; these boundary
values are such that the function In If(e iB ) I is Lebes-
gue integrable on f. The following additional property
is displayed by the classes HP, O<p < + 00, or H",: The
function I f(e,B) IF or, correspondingly, #,In I f(e ifJ ) IY'
is Lebesgue integrable on f. For bounded functions
f (z), I f (z) I :s;;;.M, one has instead of the above,
ess sup I f(e ifJ ) I ";;;'M, 0:s;;;.(J:s;;;.27T. Thus, condition (3) is
the widest sufficient condition on the average increase
of an analytic function f (z), as I z I~ 1, that ensures
the existence of angular boundary values almost every-
where on f.
It has been shown that condition (3) cannot be
substantially weakened. Thus, it was shown by A.
Zygmund that for an arbitrary increasing function 1/J(t),
 BOUNDARY PIWPERTIES OF ANALYTIC FUNCTIONS

I/I{t) / t~O as 0<t1+ 00, there exists an analytic func- it possible to extend this theorem to meromorphic func-
tion ! (z) in D such that tions of bounded characteristic (cf. Function of bounded
2w form). N.N. Luzin (1919) constructed, for any set E cf
sup f#)n+
0<,<:1 0
IJ(rei<l I)d</l < +00, of measure zero, an analytic function ! (z) such that
!(ei(J)=O everywhere on E as z~ei(J in an arbitrary
but which has no boundary values anywhere on f. manner, but !(z) is not identically equal to zero. The
Even if the maximum deepest and most general boundary uniqueness
M(r;f)=max{ I!(z) I: I z I =r} grows as slowly as theorems for meromorphic functions of general form
one pleases, there still exist analytic functions without were obtained in 1925 by Luzin and 1.1. Privalov (d.
radial boundary values. Uniqueness properties of analytic functions;
The boundary representation of functions ! (z) of Luzin - Privalov theorems).
class N(D), characterizing the functions of this class, Consider the case of an arbitrary plane domain D;
has the form for the sake of brevity, only simply-connected domains
D with a rectifiable boundary f will be discussed. The
(4)
conditions (2), (3) and (2') are equivalent to stipulating
eiO +z that the subharmonic functions I!(z) IP, In+ I!(z) I
xexp-2
1
." 0
fIn IJ(e'O) 1-o-dOX
2w

e' -z
.
and I/I{ln I!(z) I), respectively, have a harmonic
majorant in D. In such a form these conditions are fully
1 e +z iO
X exp -2
." 0
f-o-
2w

e' -z
d~(J),
where m is an integer, m =k if the point z =0 is a zero
of multiplicity k and m = -k if z =0 is a pole of multi-
plicity k; A is a real number;
I a I a -z
B I(z'a)
, JL = II--JL---JL-
00
_ (5)
JL=I aJL l-aJLz
is the Blaschke product taken over all the zeros a,,::f:O
of !(z) inside D taking into account their multiplicity;
B 2(z; b v ) is the Blaschke product of type (5) taken over
all poles bv::f:O of !(z) in D; and <1>(8) is a singular
function of bounded variation on [0,2'lT] with a deriva-
tive that vanishes almost everywhere. In (4), the last
integral is of Lebesgue - Stieltjes type, while the first is
of Lebesgue type.
It was shown by M.M. Dzhrbashyan [10] that the
theory of meromorphic functions of bounded charac-
teristic can be considerably extended. It is possible, in
fact, to introduce a family of classes of meromorphic
functions N a depending on a continuous parameter a,
-1 <a< + 00, and the classes N a are characterized by
representations yielding (4) when a=O. As a increases,
the classes N a become larger, and No becomes identi-
cal with Nevanlinna's class N.
For analytic functions !(z)EN*(D) in the represen-
tation (4) one must put B 2(z; bv)-l. For the func-
tions !(z)EHP, O<p:r;;;;, + 00, or EH.jI, in the represen-
tation (4) one has B 2 (z; b v )=I, and <1>(8) is a non-
increasing function of the indicated type. See also Cau-
chy integral.
The first results in the uniqueness problem were
obtained in 1916 by the brothers F. and M. Riesz: If a
function !(z)EHOO has radial boundary values
!(eifJ)=O on a set E Cf of positive Lebesgue measure
on f, then !(z)_O in D. The representation (4) makes
suitable, and furnish a natural definition of the classes
HP, N* and H.jI in arbitrary domains. It is known that
a rectifiable curve f has a definite tangent and normal
at almost all of its points. The inclusions (1) remain
valid, as does Fatou's theorem on the almost-
everywhere existence on f of angular boundary values
for the class N*. Here, as the bisectrix of the angular
domains ~(~, f) the normal to f at the point ~Ef is
taken. The Riesz uniqueness theorem for the class N*
is also applicable.
V.I. Smimov also introduced the frequently
employed classes EP, P >0, for the case of an arbitrary
domain D. These classes have the following definition:
!(z)EEP(D) if there exist a sequence of contours
{fj} CD, fj~f, such that
~~~ [I J(z) IP I dz I = elf,p) < +00.
The classes EP are especially convenient in the study of
problems of representation of functions in the form of
a Cauchy integral.
Of major interest is the study of the boundary pro-
perties of analytic functions realizing a conformal map-
ping. Let a function z =P(w) realize a conformal map-
ping of the unit disc I w I < I onto a domain D of the
z-plane with a rectifiable boundary f. It has been
shown, for example, that in such a case the derivative
p' (w) belongs to the Hardy class H I in the disc
I w I < 1, so that it is representable in the form (4) with
B 2 =1 and a non-increasing singular function <1>(8).
Smirnov pointed out the importance of the class S of
such domains D for which this singular function
<1>(8)_0. In 1937 M.V. Keldysh and M.A. Lavrent'ev
constructed an example of a domain with a rectifiable
boundary that is not included in the Smirnov classes
just mentioned; this renders the characterization of
domains of Smimov type even more important.

451
BOUNDARY PROPERTIES OF ANALYTIC FUNCTIONS
Numerous workers also attempt to study boundary
properties of analytic functions f (z) in several variables
Z=(Zb" . ,zn). Let D=Un={ZEC": I Zj I <I,
j=I, ... ,n} be the unit polydisc, and let Tn={zEC":
I Zj 1=1, j= 1, ... ,n} be its skeleton. The class
N*(D) of analytic functions f(z) in D is defined by the
condition
SUI>
O<""-:!r
(In+ I!(rz) I dmiz) = Clf) < +00,
which is analogous to (3), while the classes HP(D) or
H",(D) by the condition of the type (1./I(t)=e pt for the
case HP(D), p >0):
SUI> 1/;-1
O<r<::1
{ftf.{ln
f. I J(rz) I)dmn(Z)} = Clf,1/;) < +00,
where mn is the normalized Haar measure on Tn,
mn(Tn)=l. Inclusions of the type B=HooCH",CN*
are preserved. Analytic functions f(Z)EN*(D) have
r
'radial' boundary values (z) = limr ..... \-0 f (rz), Z E Tn ,
r
almost everywhere on Tn with respect to the Haar
measure mn; and In + I (z) I is summable on Tn with
respect to mn Sufficiently simple and general charac-
teristics for boundary representations and uniqueness
properties of functions f (z) in un for n > 1 have not
yet (1986) been found.
Many boundary properties may be applied to various
generalizations of analytic functions, in particular to
abstract analytic functions f: D~X, which have values
in, say, a separable locally convex topological space X
over the field C.
References
[I] GoLUBEV, V.V.: Single-valued analytic functions. Automorphic
functions, Moscow, 1961 (in Russian).
[2] GoWZIN, G.M.: Geometric theory offunctions of a complex
variable, Amer. Math. Soc., 1969 (translated from the Russian).
[3] PIuwALOW, 1.1. [1.1. PIuvALov]: Randeigenschaften analytischer
Funktionen, Deutsch. Verlag Wissenschaft., 1956 (translated
from the Russian).
[4] KHAVINSON, S.YA.: 'Analytic functions of bounded type', Itogi
Nauk. Mat. Anal. 1963 (1965), 5-80 (in Russian).
[5] NEVANLINNA, R.: Analytic functions, Springer, 1970 (translated
from the German).
[6] NOSHIRO, K.: Cluster sets, Springer, 1960.
[7] COLLINGWOOD, E.F. and LOHWATER, A.J.: The theory of clus-
ter sets, Cambridge Univ. Press, 1966.
[8] MAcLANE, G.R.: Asymptotic values of holomorphic functions,
Rice Univ. Studies, Math. Monographs, 49, Rice Univ., Hous-
ton, 1963.
[9] LOHWATER, A.: 'The boundary behaviour of analytic func-
tions', Itogi Nauk. Mat. Anal. 10 (1973),99-259 (in Russian).
[10] DZHRBASHYAN, M.M.: Integral transforms and representation of
functions in a complex domain, Moscow, 1966 (in Russian).
[II] RUDIN, W.: Function theory in polydiscs, Benjamin, 1969.
[l2] KHENKIN, G.M. and CHIRKA, E.M.: 'Boundary properties of
holomorphic functions of several complex variables', J. Soviet
Math. 5, no. 5 (1971), 612-687. (Itogi Nauk. i Tekhn. Sovrem.
Prohl. Mat. 4 (1975), 13-142)
E.D. Solomentsev
Editorial comments. SokhotskiT's theorem is better
known as the Casorati- Weierstrass theorem. In (4) the first
452
integral on the right-hand side is called the outer factor,
while the other factors on the right-hand side form together
the inner factor, the latter has absolute boundary values 1
almost-everywhere. If the outer factor in (4) equals 1, f is
called an inner function. Recent general references are
[A2], [A3].
For smooth domains D in en, n~2, there is a strong ver-
sion of Fatou's theorem for HP(D) and N'(D): Boundary
values exist almost-everywhere (or even better than almost-
everywhere) if the boundary is approached over admissible
approach regions, see [AS], [A7]. In contrast with the one-
dimensional case, zero sets of HP-functions can be essen-
tially different for different values of p, see [11], [AS]. A
characterization of zero sets of functions of Nevanlinna
class on smooth strongly pseudo-convex domains was
obtained by G.M. Khenkin and, independently, by H. Skoda,
see [AS]. A.B. Alexandrov and E. L0W proved indepen-
dently the existence of non-constant inner functions for the
unit-ball in en, see [A1], [A4]. In an other direction, HP-
spaces have been introduced for the upper half-space in
Rn , [A6].
References
[A 1] ALEKSANDROV, A.B.: 'The existence of inner functions in the
ball', Math. USSR Sb. 46 (1983), 143-159.
[A2] GARNETT, J.: Bounded analytic functions, Academic Press,
1981.
[A3] KOOSIS, P.: Introduction to Hp spaces, Cambridge Univ.
Press, 1980.
[A4] L0w, E.: 'A construction of inner functions on the unit ball in
cP', Invent. Math. 67 (1982), 223-229.
[AS] RUDIN, W.: Function theory in the unit ball in cn, Springer,
1980.
[A6] STEIN, E.M.: Singular integrals and differentiability properties
of functions, Princeton Univ. Press, 1970, Chapt. 7.
[A7] STEIN, E.M.: Boundary behaviour of holomorphic functions of
several complex variables, Princeton Univ. Press, 1972.
AMS 1980 Subject Classification: 30D40, 30055
BoUNDARY VALUE PROBLEM, COMPLEX-
VARIABLE METHODS - Methods for studying boun-
dary value problems for partial differential equations in
which one uses representations of solutions in terms of
analytic functions of a complex variable.
Given a second-order elliptic equation
~u +a(x,y)
au +b(x,y) ay
ax au +c(x,y)u = 0, (1)
where a, b, c are analytic functions of the real variables
x, yin some domain Do of the z-plane, Z =x + ry, con-
sider the following boundary value problem: Find a
solution of equation (1), regular in a simply-connected
domain S CD o, satisfying the boundary condition
=
[ aJ+k u ] ]
R(u) ~
O<;j+k<;m
[ a j k (aJ+ku
X
t) k + 1jk
ja- a
~
-X a
~
ja k = f(t)O)
where ajk(t), f(t)ECa(aS), O<a< 1, and the 1jk are
linear operators mapping Ca(aS) into Ca(aS), with
Tm-k,k completely continuous.
This problem includes as special cases the well-
 BOUNDARY VALUE PROBLEM, COMPLEX-VARIABLE METHODS

known classical boundary value problems of Dirichlet,
Neumann, Poincare, etc.
Using the formula for the general representation of
solutions (see Differential equation, partial, complex-
variable methods),
u(x, y) = Re {G(Z' zo, z, Z)41(z)-
-
Z a _ }
tP(t)-atG(t, zo, z, Z)dt ,
one reduces the problem to an equivalent boundary
value problem for analytic functions:
m
Re ~ [ak(t)41(k)(t) + Tk(41(k))] = jet),
k=O
where ak(t) are given Holder-continuous functions,
tEaS, T m is a completely-continuous operator, and the
Tk (k = 0, ... ,m - 1) are linear operators.
Suppose that the finite simply-connected domain Sis
bounded by a closed Lyapunov contour (see Lyapunov
surfaces and curves) as and that the m-th derivative,
m ;;;:'0, of the function cI>(z) (the latter is holomorphic in
S), restricted to as, is a function of class C a , 0<0:< 1.
Then, assuming that the point z =0 is in S, one can
express cI>(z) as follows:
4l(z) = !st1~:t + ie, ifm=O;
where p is the increment of the function (1/ 2'IT) arga(t)
when the contour as is described once in the positive
sense. The index of the Dirichlet problem is zero.
The homogeneous boundary value problem R(u)=O
has a finite number k;;;:.O of linearly independent solu-
tions, where k ;;;:.,,; the inhomogeneous problem (2) has
a solution if and only if
/!(t)pit)dS =0, j=I, ... ,k,
where Pi is a complete system of linearly independent
solutions of the associated homogeneous integral equa-
tion
k'(p) = A(to)P(to)-~ r B (t)1I(t)dt +T(p) = o.
'lT1 is t - to
(3) The boundary value problem (2) has a solution,
whatever the free term on the right, if and only if there
exist exactly " linearly independent solutions of the
associated homogeneous problem R(u)=O. Conse-
quently, if 1(>0 the homogeneous boundary value prob-
lem R(u)=O always has at least I( linearly independent
solutions; if ,,<0 the inhomogeneous problem (2) is not
solvable for an arbitrary free term on the right, but at
least I" I solvability conditions must be satisfied.
Necessary and sufficient conditions for the solvability
of an inhomogeneous boundary value problem may be
formulated in terms of the completeness of a certain
system of functions. The kernel and system of func-
tions may be constructed explicitly using the Riemann

[I-fr- [I-f]
function of equation (1) and the coefficients of the
<I>(z) = !s~t) 1ln dl + boundary conditions. For example, let {ud be some
complete system of solutions in the basic domain Do of
+ !s~t)dt+ie, if m;;;.I, equation (1), and let S CD o. Then a necessary and suf-
ficient condition for problem (2) to be solvable for any
where p.(t) is a real function of class C a(aS), 0<0:< 1, free right-hand side is that the system of functions
and c is a real constant; p.(t) and c are uniquely deter- {R (Uk)} be complete on the boundary.
mined by cI>(z). Very complete results have been obtained for the fol-
Substituting these expressions into the boundary con- lowing boundary value problem (the generalized
dition (3), one obtains a singular integral equation, Riemann- Hilbert problem): Find a solution of the
equivalent to problem (2), for the unknown function p.: equation

K(p.) = A(toMto)+ B(t.o) r.l!!l.J.!lL + T(p.) = j(to), ilzw+A(z)w+B(z)w = j(z), w=u+iv, (5)
'lT1 is I - to 2ilz = ax-iay,
toE as, where T is a completely-continuous operator. which is continuous in S + as and satisfies the boun-
A necessary and sufficient condition for the boun- dary condition
dary value problem (2) to be normally solvable is that
m
Re[X(z), w(z)]- au+fiv = ,(, ZEaS, (6)
aCt) = ~ikam_k(t) =1= 0, lEaS, m;;;'O. (4) where 0:, {3, y are real functions of class Ca(aS),
0<0:< 1, with 0:2 + {32 = 1. The domain S may be multi-
k=O

The Dirichlet problem (m =0) is always normally solv-

ply connected. A problem of this type may be reduced
able. (Henceforth it is assumed throughout that condi-
to an equivalent singular integral equation; this yields a
tion (4) is satisfied.)
full qualitative analysis of the boundary value problem
The index of the boundary value problem (2) is com- (6).
puted from the formula Suppose that the boundary as of S is the union of a
I( = 2(m+p), m;;;.l, finite number of simple closed curves as o, ... ,asm

453
BOUNDARY VALUE PROBLEM, COMPLEX-VARIABLE METHODS
satisfying the Lyapunov conditions. Since the forms of
the equation and the boundary condition are preserved
under conformal mapping, it may be assumed without
loss of generality that as o is the unit circle with centre
at Z =0, the latter being a point of S, while
as J, . ,aSm are circles lying outside as o.
The index of problem (6) is defmed as the integer n
equal to the increment of (l /2w)arg[a(O+i,8(O] when
r
the point describes as once in the positive sense. The
boundary condition can be reduced to the simpler form
Re[z-neic(Z)w(z)] = y, zeas,
where c(z)=Cj on asj , with co=O, while CJ, . . 'Cm
are real constants, uniquely expressible in terms of a
and ,8. The index of the adjoint problem
l
ilzw.-Aw.-Bw. = 0, zeS,
Re{<a+I/l):!w.<,J} ~ 0, ,eas, J
(7)
is calculated by the formula n' = - n + m - 1.
Problem (6) has a solution if and only if
l(a+i{I)w.yds = 0,
where w. is an arbitrary solution of the adjoint prob-
lem.
Let e and e' be the numbers of linearly independent
solutions of the homogeneous problems (6) and (7),
respectively. Then
e-e' = n-n' = 2n+l-m.
If n <0, the homogeneous problem (6) has no non-
trivial solutions. If n >m -1, it has exactly
=
e 2n + 1 - m linearly independent solutions, while the
inhomogeneous problem (6) is always solvable. If n <0,
the inhomogeneous problem (6) is solvable if and only
if
l(a+i/J)w./fds = 0, j=I,2, ... ; e'=m-2n-l,
where w.j is a complete system of solutions of the
homogeneous problem (7). If m =0 and n =0, then
e = 1 and all solutions of the homogeneous problem
problem (6) are given by
w(z) = ice""'(Z) ,
where c is a real constant and Wo a continuous function
on s+as.
The above results completely characterize the prob-
lem in the simply-connected (m =0) and multiply-
connected (n <0, n > m - 1) cases. The cases
O";;;n";;;m -1 require special examination; they have
also been worked out in considerable detail.
Boundary value problems of the type of the Poincare
problem have also been studied for equation (5).
For references see Differential equation, partial,
454
complex-variable methods.
I.N. Vekua
Editorial comments.
References
[A1] VEKUA,I.N.: Generalized analytiC functions, Pergamon, 1962
(translated from the Russian).
[A2] VElmA, N.P.: Systems of singular integral equations and
some boundary value problems, Moscow, 1970 (in Russian).
[A3] VEKUA, I.N.: New methods for solving elliptic equations,
North-Holland, 1968.
AMS 1980 Subject Classification: 35J25, 35G15,
30E25,45E05,35C15
BOUNDARY VALUE PROBLEM, ELLIPrIC EQUA-
TIONS - The problem of finding a solution u, regular
in a domain D, to an elliptic equation
n a2u n au
~ a;k--+ ~b;-+cu = j, (I)
;,k=O aX;aXk ;=0 ax;
which satisfies certain additional conditions on the
boundary r of D. Here ajb bj, C and f are given func-
tions on D.
The classical boundary value problems are special
cases of the following problem: Find a solution to
equation (1), regular in a domain D and satisfying on r
du
a-+OO
dl
=g (2)
where d / dl denotes differentiation in some direction,
and a, b and g are given continuous functions on r
with I a I + I b I >0 everywhere on r (see [1]).
Putting a =0, b = 1, one obtains the Diricblet prob-
lem; with b =0, a = lone has a problem with oblique
derivative (see Differential equation, partial, oblique
derivatives), which becomes a Neumann problem if I is
the direction of the conormal. If r = 1\ U f 2, where r I
and r 2 are disjoint open subsets of r, and fin f 2 is
either empty or an (n - 2)-dimensional manifold, with
a = 1, b =0 on rio a =0, b = 1 on r 2 , one obtains a
mixed problem.
Problem (2) has been studied for elliptic equations in
two independent variables (see [2]). Fairly complete
investigations have been made of the Dirichlet problem
for elliptic equations in any finite number of indepen-
dent variables (see [1], [3], [4]) and the problem with
oblique derivative in case the direction f is not con-
tained in a tangent plane to r at any point of r. In
that case the problem with oblique derivative is a
Fredholm problem and the solution is smooth to the
same order as the field of directions f and the function
g (see [l]). The case in which flies in a tangent plane
to r at certain points of r has been studied (see [3]).
The local properties of solutions to the problem with
oblique derivative have been investigated (see [5]). At
points where the field flies in a tangent plane to r, the
solution of the problem is less smooth than I and g.

This has been used as a basis for investigating the
problem in a generalized setting (see [7], [8)).
Consider the following boundary problem for har-
monic functions regular in the unit ball ~ CR3 :
aux + bUy +cuz = g;
let K be the set of points of the unit sphere S at which
the function w = ax + by + cz vanishes. The vector field
Pea, b, c) lies in a tangent plane to S at the points of
K. Suppose in addition that K is the union of a finite
number of disjoint curves; let K+ be the subset of K
consisting of those points at which grad w makes an
acute angie with the projection of the field P on S, and
let K- be the remaining part of K. A generalized for-
mulation of the problem is obtained when the values of
u are also prescribed on K+, whereas on K- the solu-
tion u is allowed to have integrable singularities. If
K- is empty, the solution to the generalized problem
may be made arbitrarily smooth by increasing the
smoothness of the additional data of the problem. Gen-
erally speaking, a solution to the mixed problem on the
set fo=1\ ni\ has singularities (see [1D. In order to
eliminate such singularities on fo, one must impose
additional conditions on the data (see [11)).
A large category of boundary value problems is con-
stituted by what are known as problems with free
boundaries. In these problems one must find not only a
solution of equation (1), but also the domain in which
it is regular. The boundary f of the domain is
unknown, but two boundary conditions must be satis-
fied on it. An example of this type of problem is the
problem of wave motions of an ideal fluid: Find a har-
monic function u, regular in some domain D, where
part of the boundary, fl say, is known and the normal
derivative au / an is given on f l ; the other part of the
boundary, f2' is unknown and on it one gives two
boundary conditions:
au -_ 0,
a; 2 2
ux+uy+u z
2_
- q(x,y, z),
where q >0 is a given function.
For harmonic functions of two independent vari-
ables, one uses conformal mapping (see [12], [13], [14)).
See also Differential equation, partial, free boundaries.
The following problem has been investigated: Find a
harmonic function u, regular in a domain D and satis-
fying the condition
1 gradu 12 = q,
where q>O is a given function, on the boundary f.
There is a complete solution of this problem for har-
monic functions of two independent variables (see [14)).
Given an equation Lu = f, where L is an ope~tor of
order 2m, uniformly elliptic in the closure D of a
domain D, consider the problem of determining a solu-
BOUNDARY VALUE PROBLEM. ELLIPTIC EQUATIONS
tion u, regular in D and satisfying on the boundary f
of D the conditions
B)u = ct>j, )=1, ... ,m, (3)
where BJ<x, D), j= 1, ... ,m, are differential operators
satisfying the following complementarity condition.
Let L' (x, a/ ax I, . . . ,a / aXn+l) be the principal
part of L, let B; be the principal part of Bj , n the nor-
mal to f at a point x and >..*0 an arbitrary vector
parallel to f. Let Tt (>..) denote the roots of
L' (x, >.. +m) with positive imaginary parts. The polyno-
mials B;(x, >"+m), j=l, ... ,m, as functions of T,
must be linearly independent modulo the polynomial
II;= I(T-Tt (>..)). In this case, too, the problem is nor-
mally solvable. Violation of the complementarity condi-
tion may entail an essential change in the nature of the
problem (see [17]).
Problem (2) is a special case of problem (3). For
problem (2) with a 1, the complementarity condition
is equivalent to the condition that there be no point on
the boundary of the domain at which the direction I
lies in a tangent plane to the boundary.
Another particular case of problem (3) is the boun-
dary value problem
a;u = ct>j, )=0, ... ,m-I,
an)
which is an analogue, to some extent, of the Dirichlet
problem for higher-order elliptic equations.
The boundary value problem has been studied for
the poly-harmonic equation t::.ku =0 when the boundary
of the domain consists of manifolds of different dimen-
sions (see [15)).
In investigations of boundary value problems for
non-linear equations (e.g. the Dirichlet and Neumann
problems), much importance attaches to a priori esti-
mates, various fixed-point principles (see [17], [18)) and
the generalization of Morse theory to the infinite-
dimensional case (see [19)).
References
[1) MIRANDA, c.: Partial differential equations of elliptic type,
Springer, 1970 (translated from the Italian).
[2) VEKUA, LN.: Generalized analytic functiOns, Pergamon, 1962
(translated from the Russian).
(3) BITSADZE, A.V.: Boundary value problems for secondorder ellip-
tic equations, North-Holland, 1968 (translated from the Rus-
sian).
(4) KELDYSH, M.V.: 'On the solvability and stability of the Diri-
chlet problem', Uspekhi Mat. Nauk, no. 8 (1941), 171-231 (in
Russian).
(5) HORMANDER, L.: 'Pseudo-differential operators and non-
elliptic boundary value problems', Ann. of Math 83, no. 1
(1966), 129-209.
(6) BORRELLI, R.L.: 'The singular, second order oblique derivative
problem', J. Math and Mech. 16, no. 1 (1966), 51-81.
(7) GoROV, YU.V. and KONDRAT'EV, V.A.: 'The oblique deriva-
tive problem', Mat. Sb. 78, no. 1 (1969), 148-176 (in Russian).
(8) MAz'YA, V.G.: 'The degenerate problem with oblique deriva-
tive', Mat. Sb. 87, no. 3 (1972), 417-453 (in Russian).
455
BOUNDARY VALUE PROBLEM, ELLIPTIC EQUATIONS

(9) Y ANUSHAUSKAS, A.: Dokl. Akad Nauk SSSR 164, no. 4 assumed that the problem in question is well-posed. In
(1965), 753-755.
order to justify difference methods one must investigate
(10) VISHIK, M.I. and EsKIN, G.I.: 'Sobolev-Slobodinsky spaces of
variable order with weighted norm, and their applications to the well-posedness of the difference problem and its
mixed boundary value problems', Sibirsk. Mat. Zh. 9, no. 5 convergence as the grid is refined. A difference problem
(1968),973-997 (in Russian). is said to be well-posed if it has a unique stable solu-
(11) GIRAUD, G.: Ann. Soc. Math. Polon. 12 (1934), 35-54.
(12) LAVRENT'EV, M.A.: The variational method in boundary value tion, whatever the terms on the right-hand sides of the
problems for systems of equations of elliptic type, Moscow, 1962 equations. By stability of a difference scheme one
(in Russian). means that its solution depends continuously on the
(13) NEKRASOV, A.I.: 'Exact theory of waves of stationary type on
the surface of a heavy fluid', in Collected works, Vol. I, Mos-
right-hand side, uniformly with respect to the grid spac-
cow, 1961 (in Russian). ing.
(14) GAKHOV, F.D.: Boundary value problems, Pergamon, 1966 For example, suppose one wishes to solve the Diri-
(translated from the Russian).
chlet problem for the Poisson equation in the square
(15) SoBOLEV, S.L.: Mat. Sb. 2, no. 3 (1937), 465-499.
(16) AGMON, S., DoUGLIS, A. and NIRENBERG, L.: 'Estimates near G={O<Xl,X2<1} with boundary f:
the boundary for solutions of elliptic partial differential equa-
tions satisfying general boundary conditions I, II.', Comm.
Pure Appl. Math. 17 (1964), 35-92.
(17) ScHAUDER, J.: Math. Z. 33 (1931), 602-640.
(18) LERAY, J. and ScHAUDER, J.: Ann. Sci. Ecole Norm. Sup. Ser. 3 U(X\>X2) = p.(X\>X2), (X\>X2)Ef.

51 (1934), 45-78. The domain G is replaced by a square grid Gh with

(19) PAUlS, R.S.: 'Morse theory on Hilbert manifolds', Topology 2,
spacing h, i.e. by the set of points
no. 4 (1963), 299-340. AI v h k
. . anus aus as
~
Gh = {(X\i) , x~: xV) = ih, x~) =jh ;
Editorial comments.
i,j=l, ... ,N-l}, hN=l,
References
[A1) COURANT, R. and HILBERT, D.: Methods of mathematical with boundary
physics. Partial differential equations, 2, Interscience, 1965 N-j

(translated from the German). fh = U {(O, x~ U (1, x~ U (X\i) , 0) U (xV), I)},

i,)=1
[A2] GARABEDIAN, P.R.: Partial differential equations, Wiley, 1964.
[A3] FRIEDMAN, A.: Partial differential equations, Holt, Rinehart and the derivatives figuring in the equation by differ-
and Winston, 1969.
ence quotients
AMS 1980 Subject Classification: 35JXX (Ui+ 1,) - 2Ui,} +Ui-j,)
(AjU)i.) = ux\x"i,) = h2
BOUNDARY VALUE PROBLEM, NUMERICAL
METHODS FOR PARTIAL DIFFERENTIAL EQUA- (Ui,j+ 1 -2Ui,) + U1,)_I)
(A 2 u)i,) = Ux ,x2,i,) = h2
TIONS - Approximate methods of solution which yield
the solution to the problem in the form of a numerical where Ui,j = U(X\i) , xY). The difference scheme is
table. An exact construction of a solution (in terms of
(A"y )i,) = YX\X\.I,) +Yx,x"i,) = - f,,}, (1)
explicit formulas, series, etc.) to a boundary value prob-
lem is only rarely possible. The most prevalent of the i,j=I, ... ,N-l; Yi,) 1
f,
= /-Lip
methods of approximate solution are difference
methods (see [1]); they are applied in the most general where Yi,j is the solution.
problems and are readily computerized. The essential Problem (1) has a - unique - solution for any
idea of difference methods is to replace the original inhomogeneous term f and any boundary conditions Jl
domain of variation of the independent variables by a (see [2]). Moreover, the solution of the difference
discrete set of points, a grid, and to approximate the scheme (1) converges as h ---'?O to the solution of the ori-
derivatives figuring in the equation and the boundary ginal problem in such a way that the scheme is a
conditions by difference quotients at the grid points. As second-order approximation in the maximum norm, i.e.
a result of this procedure the original problem is
replaced by a system of finitely many algebraic equa-
tions (linear or non-linear); this system is known as a where M is a constant independent of h.
difference scheme. The solution of the difference The difference scheme (1) is a system of linear alge-
scheme is taken to be an approximate solution of the braic equations, characteristically with a large number
original problem. The accuracy of the approximation of equations (to be precise: (N -1)2 equations, with
depends on the method of approximation and on the N ---'? 00 as h ---'?O) and a large number of zeros in the
fineness of the grid, i.e. on how densely the grid points matrix of the system; moreover, the system is generally
fill the original domain. In the rest of this article atten- ill-conditioned (the ratio of the minimum eigen va)ue to
tion will be restricted to linear boundary value prob- the maximum one is of the same order of magnitude as
lems for partial differential equations, and it will be h 2, h ---'?O). Such systems of equations, which arise when

456
 BOUNDARY VALUE PROBLEM, NUMERICAL METHODS FOR PARTIAL DIFFERENTIAL EQUATIONS

differential equations are replaced by difference equa- mines the stability and accuracy of the scheme. If
tions, may be solved using various efficient methods, o~0.5, the scheme is stable for any grid spacings 7 and
both direct and iterative. Direct methods yield an exact h (an absolutely-stable difference scheme). But if
solution of the difference equation after a finite number (1<0.5, the scheme is stable only if some restriction is
of arithmetic operations. This category includes various imposed on Y=7/ h 2 (a conditionally-stable difference
versions of the double-sweep method, including matrix scheme); for example, the explicit scheme 1=0) is
double-sweep, the decomposition method, the fast stable if y~ 1/4. If (1=0.5, the scheme is second-order
Fourier transform, and the method of representations accurate with respect to 7 and h; for other values of (1 it
by sums (see [1], [2], [3], [6]). A measure of the effi- is first-order accurate in 7 and second-order accurate in
ciency of direct methods is the order of magnitude of h. The difference problem (3) is solved 'in layers'. The
the number of operations as h ~O. Thus, solving prob- n-th layer is the set of all points of the grid Gh X w,. for
lem (1) by matrix double-sweep requires O(h -4) opera- some fixed n. The values of Y?,j on the zero-th layer
tions, whereas the fast Fourier transform method needs (n = 0) are known from the initial conditions. If the
- for the same problem - O(h -210g h -\) operations. values on the n-th layer are known, then the values on
The most widely used iterative methods for the solution the next layer Y;,j = y7,f \ are determined from the sys-
of difference problems are Richardson's method with tem of equations
Chebyshev interpolation points, the alternating-
triangular iterative method, and various alternating-
direction methods (see [2]). The efficiency of iterative
Yi'J-o~lll!}')i.j = fl.j, i,j=I, ... ,N-I, 1 (4)
YI.jl - 0,
r,
methods is measured by the order of magnitude of the
minimum number of iterations, no(), necessary to where
decrease the error of the initial approximation by a fac- fl.j = Y2j+(1-o)T(1l1!}')2 r
tor 1/ . For example, in solving problem (1) by The solution of problem (4) may be found by any of
Richardson's method, no() is of the order the methods for solving the stationary problem (1).
h -\ log(l / ), while the alternating-direction method However, there are more economical algorithms for the
with optimally-selected iteration parameters has solution of multi-dimensional non-stationary boundary
no()= O(logh -\ 10g(1 / . Iterative methods are more value problems, namely alternating-direction methods
universal and simpler to implement than direct (see [1] - [5]), which enables one to reduce the solution
methods, and consequently have been widely utilized to of a multi-dimensional problem to that of a sequence of
solve difference problems. one-dimensional problems. Thus, the heat equation
For example, consider solving the first boundary may be solved using the following alternating-direction
value problem for the heat equation: scheme:
au au au 2 2 n+I/2 n
-a
(
= -aXI
2 +-,
ax~
(>0, (Xj, X 2)EG Yi.j
0.57
-Y,.j = A n+I/2+A n
IY,.j zY,.j'
n+1 _ n+I/2
U(Xj,X2,0) = UO(XI,X2); (2) Y i.j Y i.j = A n+ I /2 + A n+ I
0.57 IY,.j zY,.j .

This scheme is absolutely stable, second-order accurate

To solve this problem, one prescribes a time grid with
and is solved by successive inversion of one-
spacing 7>0:
dimensional difference operators.
WT = {tn=n7: n=O, I, ... } Even when the difference problem is solved exactly,
and the grid Gh (see above) for the space variables. Let the solution may well be different - not only quantita-
u7,j = u(x\i), x~), tn). The derivative u; is approximated tively but also qualitatively - from the solution of the
by the quotient U n + 1 -un.
original differential problem. The discrepancy is
I.j I.j
Ul.I.j = particularly prominent in dealing with equations the
7
coefficients - or solutions - of which involve singu-
and the Laplacian Il by the difference operator Ilh . larities (thus, in computing discontinuous solutions of
The original equation (2) is replaced by a correspond- the equations of gas dynamics one usually encounters
ing difference scheme zones in which the discontinuities are strongly 'spread
ouC). Thus, direct approximation of a differential prob-
Yt. i.j = o( Ill!}' )7,/ I + (I - 0)( Ill!}' )~j ; lem by a difference problem, when the derivatives are
y?,j = uo(x\i),xY, i,j=I, ... ,N-I; (3) replaced - with a considerable degree of arbitrariness
Y2jl
f",xw,.
= o. - by difference quotients, does not always yield a
good difference scheme. Various principles have been
The parameter (1 appearing ill these equations deter- developed for the construction of difference schemes of

457
BOUNDARY VALUE PROBLEM, NUMERICAL METHODS FOR PARTIAL DIFFERENTIAL EQUATIONS

good quality. Among the successful methods are the lying in a given subset D of the space D(J, Rn) of func-
balance method and the method of variational- tions depending on t that are absolutely continuous on
functional approximation (see [1] - [3]). The difference J and that assume values in Rn:
schemes obtained in these methods correctly reflect the x() E D. (2)
integral laws of conservation which are valid for the
It is assumed that f(t, x) is a function defined on
original equations, and ensure that the relevant differ-
J X Rn with values in Rn and satisfying the
ence operators have fixed signs. In the theory of homo-
Caratheodory conditions; J is an interval on the real
geneous difference schemes (see [7]) one considers ques-
line R.
tions of constructing difference schemes for equations
1) The boundary value problem (1), (2) is said to be
with variable (including discontinuous) coefficients and
investigating their convergence.
linear if
jet, x) = A (t)x +b(t),
References where the functions A(t) and b(t) are summable on
[I] TncHONOV, A.N. and SAMARSKIi, A.A.: The equations of
mathematical physics, Moscow, 1977 (in Russian).
every compact interval in J and the set D is a linear
[2] SAMARSKIi, A.A. and MARCHUK, G.I.: The theory of difference manifold in D(J, Rn). In particular, one might have
schemes, Moscow, 1983 (in Russian).
[3] MARCHUK, G.I.: Methods of computational mathematics, Mos- J = [to. td,
cow, 1977 (in Russian).
[4] YANENKO, N.N.: Fractional step methods for solution of multi-
dimensional problems in mathematical physics, Novosibirsk,
1967 (in Russian).
D ~ {xo E D(J, R'), [!d<l>(t)]X(t) ~ j, 0

[5] D'YAKONOV, E.G.: Difference methods for solving boundary

value problems, 1-2, Moscow, 1971-1972 (in Russian).
[6] POLOZHll, G.N.: Numerical solution of two- and three-
dimensional boundary value problems of mathematical physics
and functiOns of a discrete argument, Kiev, 1962 (in Russian).
[7] TncHONOV, A.N. and SAMARSKIi, A.A.: 'On homogeneous
difference schemes', Zh. Vychisl. Mat. i Mat. Fiz. 1, no. 1
(1961), 5-63 (in Russian).
A.V. Gulin
Editorial comments. In Western literature, boundary
value problems are usually called initial-boundary value
problems if the time variable is involved (e.g. problem (2) of
this article). The literature on the numerical solution of par-
tial differential equations is overwhelming; only a few good
text books are listed below. References [A5] and [A8] only
deal with initial-boundary value problems.
References
[A1] FORSYTHE, G.E. and WASOW, W.R.: Finite difference methods
for partial differential equations, Wiley, 1960.
[A2] GARABEDIAN, P.R.: Partial differential equations, Wiley, 1964.
[A3] GLADWELL, I. and WAIT, R. (EDS.): A survey of numerical
methods for partial differential equations, Clarendon Press,
1979.
[A4] MITCHELL, A.R. and GRIFFITHS, W.F.: The finite difference
method in partial differential equations, Wiley, 1980.
[A5] RICHTMYER, R.D. and MORTON, K.W.: Difference methods
for initial value problems, Wiley, 1967.
[A6] SAMARSKIT, A.A.: Theorie der Differenzverfahren, Akad. Ver-
lagsgesellsch. Geest u. Portig K.-G., Leipzig, 1984.
[A7] SMITH, G.D.: Numerical solution of partial differential equa-
tions, Oxford Univ. Press, 1977.
[A8] YANENKO, N.N.: The method of fractional steps; the solution
of problems of mathematical phYSICS in several variables,
Springer, 1971 (translated from the Russian).
AMS 1980 Subject Classification: 65MXX, 65NXX
BOUNDARY VALUE PROBLEM, ORDINARY DIF-
FERENTIAL EQUATIONS - The problem of finding a
solution to an equation
~ = jet, x), tEJ, xERn, (1)
where cI>(t) is a function of bounded variation. A linear
boundary value problem gives rise to a linear operator
Lx(t) = x' -A (t)x, x(-)ED,
the eigen values of which are precisely those values of
the parameter A for which the homogeneous boundary
value problem
x' -A (t)x = Ax, x(')ED,
has non-trivial solutions. These non-trivial solutions are
the eigen functions of the operator L. If the inverse
operator L -\ exists and has an integral representation
x(t) = L -1 b(t) j G(t, s )b(s) tis, t EJ,
J
then G(t, s) is called a Green function.
2) Let J=(-oo, 00), letf(t, x) be almost-periodic in
t uniformly in x on every compact subset of Rn and let
D be the set of almost-periodic functions in t that are
absolutely continuous on J. Then problem (1), (2) is
known as the problem of almost-periodic solutions.
3) In control theory one considers boundary value
problems with a functional parameter: a control. For
example, consider the equation
~~ = j(t, x, u), tEJ=[to, td, xERn, (3)
with set of admissible controls U and two sets
Mo, M\ eRn. Let D be the set of absolutely continuous
functions in t such that X(tO)EMO, x(t\)EM\. The
boundary value problem is to find a pair (xoo, uoO)
such that UOOE U and the solution xo(t) of equation
(3) at u=uo(t) satisfies the condition xoOED.
4) There is a wide range of diverse necessary and suf-
ficient conditions for the existence and uniqueness of
solutions to various boundary value problems, and of

458
 BOUNDARY VALUE PROBLEM, PARTIAL DIFFERENTIAL EQUATIONS

methods for constructing an approximate solution (see [A2] INCE, E.L.: Ordinary differential equations, Dover, reprint,
1956.
[4] - [7]). For example, consider the problem [A3] JACKSON, L.K.: 'Boundary value problems for ordinary dif-
ferential equations', in J.K. Hale (ed.): Studies in ordinary dif-
x' = A (t)x +f(t, x), }
ferential equations, Math. Assoc. Amer., 1977, pp. 93-127.
I, (4)
j[dq,(t)]x(t) = 0, AMS 1980 Subject Classification: 34BXX
10
in which
IIf(t, x) II ~a+bllxlla BOUNDARY VALUE PROBLEM, PARTIAL DIF-
FERENTIAL EQUATIONS - The problem of determin-
for certain constants a >0, b >0, a~O. Suppose that
ing in some region D with points x=(xJ, ... ,xn ) a
the homogeneous problem
solution u(x) to an equation

I,

x' = A(t)x, j[dq,(t)]x(t) = (5) (LuXx) = f(x), XED, (1)

10
which satisfies certain boundary conditions on the boun-
is regular, i.e. its only solution is the trivial one. Then dary S of D (or on a part of it):
problem (4) has at least one solution, provided either
a< 1, or a~ 1 and b is sufficiently small. It is fairly (Bu)(y) = </>(y), YES. (2)
complicated to determine whether problem (5) is regu- As a rule, the boundary conditions relate the boundary
lar. However, the linear (scalar) boundary value prob- values of the solution to its derivatives up to a certain
lem order, i.e. B is a differential operator. However, boun-
XU +q(t)x' +p(t)x = 0, x(to) = 0, x(t]) = 0, dary conditions of other types also occur.
Given a differential equation, the question whether a
for example, is regular if whenever I q(t) I :;;;;2m there
specific boundary value problem should be studied is
exists a k ER such that
frequently settled by the concept of what is known as
I,
well-posedness. Namely, a boundary value problem is
jfp(t)-k]+ dt < 2[F(k, m)-mJ,
10
well-posed if it is solvable, and if its solution is unique
and depends continuously on the data of the problem.
where
Different types of differential equations require dif-
_~ (t]-to)~ ferent well-posed boundary value problems; and con-
V k-m" cotg 2 ' versely, well-posed boundary value problems may
sometimes serve as a basis for the classification of types
m2 <b;;;;;m 2+ .,T2 ,
(t] -tof of differential equations.
F(k, m) = A boundary value problem is said to be linear if the
operators L and B are linear, and homogeneous if f and
cJ> in (l), (2) vanish. A linear boundary value problem
is said to be Noetherian if: a) the homogeneous prob-
lem has a finite number k of linearly independent solu-
tions; b) the inhomogeneous problem is solvable if and
References
[I] HARTMAN, P.: Ordinary differential equations, Birkhauser, 1982. only if f and cJ> satisfy I linearly independent ortho-
[2] KRAsNOSEL'sKIi, M.A., BURD, V.SH. and KOLESOV, Yu.S.: gonality conditions; and c) on the assumption that the
Nonlinear almost-periodic oscillations, Wiley, 1973 (translated problem is uniquely solvable, the solution depends con-
from the Russian).
[3] PONTRYAGIN, L.S., ET. AL.: The mathematical theory of optimal
tinuously on f and cJ>.
processes, Interscience, 1962 (translated from the Russian). If k = I, the problem is called a Fredholm problem.
[4] KRAsoVsKll, N.N.: Theory of control of motion. Linear systems, The difference k -I defines the index of the problem.
Moscow, 1968 (in Russian).
A broad range of boundary value problems for linear
[5] ZUBOV, V.l.: Lectures in control theory, Moscow, 1975 (in Rus-
sian). second-order differential equations
[6] KAMKE, E.: Differentialgleichungen: LOsungsmethoden und
LOsungen, I, Chelsea, reprint, 1971.
[7] KIGURADZE, I.T.: Some singular boundary value problems for
ordinary differential equations, Tbilis~ 1975 (in Russian).
= f(x)
Yu. V. Komlenko
E.L. Tonkov falls into the category of the Poincare problem. In this
Editorial comments. type of problem the boundary conditions are prescribed
References
on the entire boundary, which is assumed to be an
[A 1] BRAUN, M.: Differential equations and their applications, (n - 1)-dimensional manifold, and the boundary opera-
Springer, 1975. tor B in (2) has the form

459
BOUNDARY VALUE PROBLEM, PARTIAL DIFFERENTIAL EQUATIONS
n au
(Bu)(y) = ;~;(y)~+p(y)u = CP(Y),
The Poincare problem for bounded regions D with
sufficiently-smooth boundaries have been thoroughly
investigated in the case of uniformly-elliptic operators
L.
On the assumption that the coefficients of the opera-
tors L and B in (3), (4) are sufficiently smooth, and
that the boundary of D is sufficiently smooth, the Poin-
care problem is Noetherian if n =2, >0, and is ar
Fredholm problem if n > 2, r
>0 and the vector
p=(pJ, ... ,Pn) is not tangent to S. Investigations of
the Poincare problem in the two-dimensional case make
extensive use of the theory of functions of a complex
variable (see Boundary value problem, complex-variable
methods).
Boundary conditions for the general elliptic equation
of order 2m
alai u
Lu = ~ a,,(x)-- = I(x), xED,
1"1,.;2m ax a
may be prescribed in terms of linear differential opera-
tors
alai u
Bju = ~ bja(Y)-a- = <l>j(Y), YES,
lal";m, ax
of order mj <2m, with coefficients defined on the boun-
dary S of D. Here a=(aj, ... ,an), a; is a non-negative
integer, Ia 1= L7=ja; and
alai alai
ax a ax~l ... a~:
If the operators Land Bj satisfy the so-called com-
plementarity conditions, then bounds for the derivatives
of order 2m of the function u satisfying the boundary
conditions can be formulated (relative to an appropri-
ate norm) in terms of the norm of fin (5) and suitable
norms of the boundary functions </>j in (6). Boundary
value problems of this type are said to be coercive.
Another type of boundary value problems are known
as mixed problems (cf. Mixed problem), in which dif-
ferent boundary conditions are prescribed on adjacent
sections of the boundary.
Characteristic for boundary value problems of dif-
ferential equations that are uniformly elliptic in D is
that the boundary conditions are prescribed on the
entire boundary. If the operator L in (3) is elliptic in
the interior of the region and parabolically degenerates
on a section S 0 ~ S, then, depending on the type of
degeneracy, So can be eliminated from the specification
of the boundary conditions. See also Boundary value
problem, elliptic equations.
If equation (3) is not elliptic, part of the boundary
may usually be eliminated from the specification of the
boundary conditions. For example, given the simplest
460
equation of parabolic type, the heat equation
YES. (4)
Lu = a2 u -~ = 0,
ax 2 at
the Dirichlet problem in the domain D bounded by the
straight lines x = + 1, t = + I amounts to the specifica-
tion of boundary values u on the intervals
{t=+l, -1";x";+I} and {x=+l, -1..;t";1}.
Boundary value problems for equations of mixed
elliptic-hyperbolic type are formulated in a special
manner.
Boundary value problems occupy an important posi-
tion in the theory of analytic functions. Let S be a
piecewise-smooth curve in the plane, i.e. a union of a
finite number of simple oriented arcs. The linear
conjugation problem consists in determining a function
</>(z), analytic outside S, having limit values </>(t), tES,
from both sides of S and satisfying the boundary con-
dition
(5) where G(t) and get) are given functions. The solution
to this problem can be specified in explicit form with
the help of representations of analytic functions by
Cauchy-type integrals, subject to certain additional
. assumptions concerning the functions G, g and the
l';;;;j';;;;m, (6)
curve S. See Boundary value problems of analytic func-
tion theory.
Various methods are available for the investigation of
boundary value problems. The Schwarz alternating
method, Poincare's related balayage method and the
Perron method are based on an application of the max-
imum principle. A solution using integral equations is
based on various integral representations of the solu-
tions. Among functional methods is the investigation of
boundary value problems using a priori estimates. The
theory of distributions is widely applied. In practical
applications, various finite-difference methods are quite
frequently used.
References
[lJ BERS, L., JOHN, F. and SCHECHTER, M.: Partial differential
equations, Interscience, 1964.
[2J BITSADZE, A.V.: Boundary value problems for secondorder ellip-
tic equations, North-Holland, 1968 (translated from the Rus-
sian).
[3J COURANT, R. and HILBERT, D.: Methods of mathematical phy-
sics. Partial differential equations, 2, Interscience, 1965
(translated from the German).
[4J LADYZHENSKAYA, O.A.: Boundary value problems ofmathemati-
cal physics, Moscow, 1973 (in Russian).
[5] MIRANDA, c.: Partial differential equations of elliptic type,
Springer, 1970 (translated from the Italian).
[6] MUSKHELISHVILI, N.!.: Singular integral equations, Wolters-
Noordhoff, 1972 (translated from the Russian).
[7] TICHONOFF, A.N. [A.N. TIKHONOV] and SAMARSKIi, A.A.: Dif
ferentialgleichungen der mathematischen PhYSik, Deutsch. Ver-
lag Wissenschaft., 1959 (translated from the Russian).
A.P. Soldatov
Editorial comments. Useful additional references are
 BOUNDARY VALUE PROBLEMS OF ANALYTIC FUNCTION THEORY
given below.
References
[A 1) COURANT, R. and HILBERT, D.: Methods of mathematical
physics, 1-2, Interscience, 1953-1962 (translated from the
German).
[A2) FRIEDMAN, A.: Partial differential equations, Holt, 1969.
[A3) GARABEDIAN, P.R.: Partial differential equations, Wiley, 1964.
[A4) HORMANDER, L.: Linear partial differential operators,
Springer, 1963.
[A5) LADYZHENSKAYA, O.A.: Linear and quasilinear elliptic equa-
tions, Acad. Press, 1968.
[A6) LIONS, 1.L. and MAGENES, E.: Non-homogenous boundary
value problems and applications, 1-3, Springer, 1972
(translated from the French).
[A7) PROTTER, M.H. and WEINBERGER, H.F.: Maximum principles
in differential equations, Prentice Hall, 1967.
[A8) FORSYTHE, G.E. and WASOW, W.R.: Finite-difference methods
for partial differential equations, Wiley, 1960.
AMS 1980 Subject Classification: 35F15, 35F30,
35G15,35G30
BOUNDARY VALUE PROBLEMS IN POTENTIAL
THEORY - Fundamental problems in both classical
and abstract potential theory. The classical Newton and
logarithmic potentials satisfy certain partial differential
equations of elliptic type: the Laplace equation in
regions free of the masses generating the potentials, and
the Poisson equation in regions occupied by the masses.
It follows, therefore, that the boundary value problems
of potential theory are primarily boundary value prob-
lems for elliptic equations and systems (cf. Boundary
value problem, eUiptic equations).
I) The Dirichlet problem, or first boundary value prob-
lem. Here the problem is to find a potential u(x) in
some domain D, given its continuous restriction
u(x) = f(x), x E r, to the boundary aD = r of the
domain on the assumption that the mass distribution in
the interior of D is known. This is a fundamental prob-
lem in potential theory.
2) The Neumann problem, or second boundary value
problem. Here the problem is to find a potential in D,
given the continuous restriction of its normal derivative
~ = Ijl(x)
an
to r.
3) The mixed problem, or third boundary value prob-
lem. Here the known data on r
constitute a linear
combination
au(x)
an +a(x)u(x) = Ijl(x), xEf. (*)
4) The problem with oblique derivative arises when
the normal derivative au(x) / an in condition (*) is
replaced by the derivative au(x) / al with respect to an
arbitrary direction l=l(x), x Er.
In addition to these general problems, the following
specific problems have also arisen in potential theory.
5) The Robin problem. Here it is required to find a
mass distribution on r with a constant potential in the
interior of D. This problem appears in electrostatics,
when it is required to determine an equilibrium charge
distribution on a conductor r which has no effect any-
where inside D.
6) The balayage method, which, in its simplest
presentation dating back to H. Poincare, amounts to
finding a 'swept-out' mass distribution on r the poten-
tial of which in the complementary domain CD coin-
cides with the potential of given masses in the interior
of D. The two last problems are particularly important
in abstract potential theory (see Potential theory,
abstract). See also Bessel potential; Non-linear poten-
tial; Riesz potential.
References
[1) GONTER, N.M.: Potential theory and its applications to basic
problems of mathematical phYSiCS, F. Ungar, New-York, 1967
(translated from the French).
[3) BRELOT, M.: Elements de la tMorie classique du potential, Sor-
bonne Uniy. Centre Doc. Uniy., Paris, 1959.
[4) CONSTANTINESCU, C. and CORNEA, A.: Potential theory on har-
monic spaces, Springer, 1972.
E.D. Solomentsev
Editorial comments. Well-known additional references
are [A 1], [A2], [A3]. Boundary value problems for parabolic
equations can be considered in the same way as those in
potential theory, cf. [A3]. The relation between boundary
value problems and the theory of Markov processes is con-
sidered in [A3].
References
[A 1) KELLOGG, O.D.: Foundations of potential theory, F. Ungar,
1929. Re-issue: Springer, 1967.
[A2) HELMS, L.L.: Introduction to potential theory, Wiley (Intersci-
ence), 1969 (translated from the German).
[A3) DOOB, 1.L.: Classical potential theory and its probabilistic
counterpart, Springer, 1983.
AMS 1980 Subject Classification: 31 A25, 31820
BOUNDARY VALUE PROBLEMS OF ANALYTIC
FUNCflON THEORY - Problems of finding an analytic
function in a certain domain from a given relation
between the boundary values of its real and its ima-
ginary part. This problem was first posed by B.
Riemann in 1857 [l]. D. Hilbert [2] studied the boun-
dary value problem formulated as follows (the
Riemann- Hilbert problem): To find the function
P(z) = u + iv that is analytic in a simply-connected
domain S + bounded by a contour L and that is con-
tinuous in S + U L, from the boundary condition
Re(a+ib)4> = au-bv = c, (1)
where a, b and c are given real continuous functions on
L. Hilbert initially reduced this problem to a singular
integral equation in order to give an example of the
application of such an equation.
The problem (1) may be reduced to a successive solu-
tion of two Dirichlet problems. A complete study of the
461
BOUNDARY VALUE PROBLEMS OF ANALYTIC FUNCTION THEORY

problem by this method may be found in [3]. Np. _A(to)lJ(to)+ [N(to,t)lJ(t)ds = (5)
The problem arrived at by H. Poincare [4] in
developing the mathematical theory of tide resembles = !(to)-c(J(to),
problem (1). Poincare's problem consists in determining
where p. is the unknown real-valued function of class H,
a harmonic function u (x, y) in a domain S + from the
c is an unknown real constant, and
following condition on the boundary L of this domain:
au au N(t t) = K(to, t) .
A(s)-a;;-+B(s)a.;-+C(s)u = !(s), (2) 0, t - to

where A (s), B (s), C (s) and f (s) are real functions
given on L, s is the arc abscissa and n is the normal to
L.
The generalized Riemann- Hilberl- Poincare problem
is the following linear boundary value problem: To find
an analytic function I>(z) in S + from the boundary
condition Re{X<Il} = !(to), toES+,
where A is an integro-differential operator defined by
the formula
X<Il = j~ {alto)<Il(f) (to) + [hlt o, t)~(f)(t)d+
where ao(to), ... ,am(/) are (usually complex-valued)
functions of class H defined on L (i.e. satisfying a
Holder condition), f (I) is a given real-valued function
of class H and hilo, I) are (usually complex-valued)
functions on L of the form
h(t t) = hJ(to, t) O";;;;a< 1,
J 0, It-tola'
where hJ(to, I) are functions of class H in both vari-
The functions A(to), K(/o, I) and a(/o) are expressed in
terms of aj(/) and hilo, 1),j=O, ... ,m.
Let k and k' be the numbers of linearly independent
solutions of the homogeneous integral equation N p. =0
corresponding to (5) and of the homogeneous integral
equation
(3) N~ - A(to)p(to)+ jN(t, to)p(t)ds = 0, (6)
L
associated with it. The numbers k and k' are connected
with the index " of the Riemann - Hilbert - Poincare
(4)
problem by the equality
K = k-k'.
Of special interest is the case when the problem is
solvable whatever the right-hand side f (10). In order
for the Riemann - Hilbert - Poincare problem to be
solvable whatever the right-hand side f(to), a necessary
and sufficient condition is k' =0 or k' = I, and in the
latter case the solution v(/) of equation (6) must satisfy
the condition
j p(t)(J(t) ds =1= 0;
L

ables. The expression 1>(;)(/0) on the right-hand side of in both cases ,,;;;;.0 and the homogeneous problem
(4) is understood to mean the boundary value on L Re{;\I>} =0 has exactly K+ I linearly independent solu-
from inside the domain S + of the j-th order derivative tions. If a(/)=O, then the Riemann-Hilbert-Poincare
of I>(z). problem is solvable for any right-hand side if and only
A special case of the Riemann - Hilbert - Poincare if k' =0.
problem, in the case when m =0, hilo,/)=O, is the As regards the Riemann - Hilbert problem, the fol-
Riemann - Hilbert problem; Poincare's problem is also lowing statements are valid: I) If ,,;;;;'0, then the inho-
a special formulation of the same problem. Many mogeneous problem (1) is solvable whatever its right-
important boundary value problems - such as boun- hand side; and 2) if ,,< -2, then the problem has a
dary value problems for partial differential equations of solution if and only if
elliptic type with two independent variables - may be 2"
reduced to the Riemann - Hilbert - Poincare problem. je i (kh/2}<t>Q(<p)c(<p)d<p = 0, k=1, ... ,-K-1,
o
The Riemann - Hilbert - Poincare problem was also
where
posed for am (to)=i=0, 10 EL, and was solved by I.N.
Vekua [3].
Q(<P) = V 1 {1 2"
exp --4 jO(<PI)cotg-2-d<P1 ,
<PI - <p }
An important role in the theory of boundary value a 2 (<p) + b 2 (<P) 7T 0

problems is played by the concept of the index of Ihe

problem - an integer defined by the formula O(t)=arg[-t-':~~~]' a 2 +b 2 =1=O.
K = 2(m +n),
The Riemann - Hilbert problem is closely connected
where 27Tn is the increment of argam(t) under one com- with the so-called problem of linear conjugation. Let L
plete traversal of the contour L in the direction leaving be a smooth or a piecewise-smooth curve consisting of
the domain S + at the left. closed contours enclosing some domain S + of the com-
The Riemann - Hilbert - Poincare problem IS plex plane z = x + ry, which remains on the left during
reduced to a singular integral equation of the form traversal of L, and let the complement of S + U L in

462

the z-plane be denoted by S -. Let a function <I>(z) be
given, and let it be continuous in a neighbourhood of
the curve L, everywhere except perhaps on L itself. One
says that the function <I>(z) is continuously extendable
to a point tEL from the left (or from the right) if <I>(z)
tends to a definite limit <1>+ (t) (or <1>- (t as z tends to
t along an arbitrary path, while remaining to the left
(or to the right) of L.
The function <I>(z) is said to be piecewise analytic
with jump curve L if it is analytic in S + and S - and
is continuously extendable to any point tEL both from
the left and from the right.
The linear conjugation problem consists of determin-
ing a piecewise-analytic function <I>(z) with jump curve
L, having finite order at infinity, from the boundary
condition
where G(t) and get) are functions of class H given on
L. On the assumption that G(t):=O everywhere on L,
the integer
is called the index of the linear conjugation problem.
If <I>(Z)=(<I>l, ... ,<I>n) is a piecewise-analytic vector,
G(t) is a square (n Xn)-matrix and g(t)=(gl' ... ,gn)
is a vector, and if also det G(t)==O, then the integer
1 method of obtaining a new solution ~=I/;Kcf>-1, ~=l/;tV1
K = 27T [argdetG(t)lL
is called the total index of the linear conjugation problem.
The concepts of the index and the total index play an
important role in the theory of the linear conjugation
problem [5], [6], [7].
The theory of one-dimensional singular integral
equations of the form (5) was constructed on the basis
of the theory of the linear conjugation problem.
References
[I] RIEMANN, B.: Gesammelte mathematische Werke - Nachtrdge,
Teubner, 1892-1902 (translated from the German).
[2] HILBERT, D.: Grnndzuge einer allgemeinen Theorie der Iinearen
Integralgleichungen, Chelsea, reprint, 1953.
[3] VEKUA, LN.: Trudy Tbil. Mat. Inst. Akad. Nauk GrnzSSR 11
(1942), 109-139.
[4] POINCARE, H.: Lefons de mecanique celeste, 3, Paris, 1910.
[5] MUSKHELISHVILI, N.I.: Singular integral equations, Wolters-
N oordhoff, 1972 (translated from the Russian).
[6] GAKHOV, F.D.: Boundary value problems, Pergamon, 1966
(translated from the Russian).
[7] KHvEDELIDZE, B.V.: Trudy Tbil. Mat. Inst. Akad. Nauk
GruzSSR 23 (1956),3-158.
A. V Bitsadze
Editorial comments. The problem discussed in the arti-
cle is also known as the barrier problem. For applications in
mathematical physics, see [A6], [A?], [A9] , and the refer-
ences given there. An important contribution to the theory
(matrix case) was given in [A5]. Other relevant publications
are [A1], [A2], [A3], [A4] and [AS]. The method proposed in
[A 1] employs the state space approach from systems
theory.
BOUNDARY VALUE PROBLEMS OF ANALYTIC FUNCTION THEORY
Note that the various names given to various variants of
these problems are by no means fixed. Thus, what is called
the linear conjugation problem above is also often known as
the Riemann-Hilbert problem [A9]. This version, especially
the matrix case where g(t), 111+ (t), 111- (t) are all (invertible)
matrix-valued functions, is of great importance in the theory
of completely-integrable systems. Indeed, consider an over-
determined system of linear partial differential equations (cf.
[A6] for more detail)
cf>x = ucf>, cf>t = vcf>, (A1)
where u, v are to be thought of as rational functions in a
complex parameter ;\. with coefficients depending on x, t but
with the pole structure independent of x, t, e.g.
U1 Un
U = Uo+--,
a1 -/\
+ ... +--,
a -/\
'
n
with the a, constants and the u, functions of x, t only. An
invertible matrix solution cf> of (A 1) exists if and only if the
corresponding u, v satisfy
ut-vx+[u, v] = 0, (A2)
a so-called Zakharov- Shabat system. Many integrable sys-
tems can be put in this form. Now let cf> solve (A 1) and take
a function g(;\.) on a contour r in the A-plane. Solve the
(x, t)-family of matrix Riemann - Hilbert problems
111- =cf>Qc/> -1111+. Then 1/;= (111- )-1 cf> also solves (A 1) and this
leads to an action of the group of invertible matrix-valued
functions in ;\. on the space of solutions of (A2). This
from an old one u, v and a function g(A) is k~o::vn as the
Zakharov- Shabat dressing method. (u, v)r->( U, v) is also
sometimes known as a Riemann- Hilbert transformation. In
the case of Einstein's field equations (axisymmetric solu-
tions) a similar technique goes by the names of
Hauser- Ernst or Kinners/ey- Chitre transformations, and in
that case (a subgroup of) the group involved is known as
the Geroch group [A 10]. The Riemann monodromy problem
asks for n multi-valued functions yeA) =(Y1 (A), ... ,Yn(A
regular everywhere but in ;\.=a1,.. ,an, oo=ao, such that
analytic continuation around a contour containing exactly
one of these points changes Y(Af into V;Y(A)T, i=O, ... ,n.
This problem reduces to the Riemann- Hilbert problem by
taking a contour through ao, ... ,an and a suitable step
function on it. The Riemann monodromy problem was
essentially solved by J. Plemelj [A11], GD. Birkhof [A12],
and I.A. Lappo-Danilevsky [A13].
References
[A1] BART, H., GOHBERG, Land KAASHOEK, M.A.: 'Fredholm
theory of Wiener-Hopf equations in terms of realization of
their symbols', Integral Equations and Operator Theory 8
(1985), 590-613.
[A2] BOUTET DE MONVEL, L., A.O. (EDS.): 'Mathematique et phy-
sique', in Sem. ENS 1979-1982, Birkhauser, 1983.
[A3] CLANCEY, K. and GoHBERG, I.: 'Factorization of matrix
functions and singular integral operators', in Operator
Theory: Advances and Applications, Vol. 3, Birkhauser,
1981.
[A4] GoHBERG, I.e. and FELDMAN, LA.: Convolution equations
and projection methods for their solution, Transl. Math.
Monographs, 41, Amer. Math. Soc., 1974.
463
BOUNDARY VALUE PROBLEMS OF ANALYTIC FUNCTION THEORY
[A5] GoHBERG, I.e. and KRETN, M.G.: 'Systems of integral equa-
tions on a half line with kernels depending on the difference
of arguments', Amer. Math. Soc. Transl. (2) (1960),217-
287. (Uspekhi Mat. Nauk 13, no. 2 (80) (1958),3-72)
[A6] ZAKHAROV, V.E. and MANAKOV, S.V.: 'Soliton theory', in
J.M. Khalatnikov (ed.): Physics reviews, Vol. 1, Harwood
Acad. Publ., 1979, pp. 133-190.
[A7] MEISTER, E.: Randwertaufgaben der Funktionentheorie,
Teubner, 1983.
[A8] RODIN, Y.L.: The Riemann boundary value problem on
Riemannian manifolds, Reidel, 1987.
[A9] CHUDNOVSKY, D. and CHuDNOVSKY, G., (EDS.): The
Riemann problem, complete integrability and arithmetic
applications, Lecture Notes in Math., 925, Springer, 1982.
[A 10] HOENSELAERS, e. and DIETz, W.: Solutions of Einstein's
equations: techniques and results, Lecture Notes in Phys.,
265, Springer, 1984.
[A11] PLEMED, J.: Problems in the sense of Riemann and Klein,
Interscience, 1 ~64.
[A12A] BIRKHOFF, G.D.: 'Singular pOints of ordinary linear differen-
tial equations', Trans. Amer. Math. Soc. 10 (1909), 436-
470.
[A 128] BIRKHOFF, G.D.: 'A simplified treatment of the regular
singular point', Trans. Amer. Math. Soc. 11 (1910), 199-
202.
[A 13] LAPpo-DANILEVSKY, LA.: Memoires sur la tMorie des
systemes des equations differentielles lineaires, Chelsea,
reprint, 1953.
AMS 1980 Subject Classification: 30E25, 45EXX,
45F15
BOUNDARY VARIATION, METHOD OF - A
method of studying univalent functions (cf. Univalent
function), based on the study of the variation of a func-
tion w = f(z) that is univalent in a domain of the z-
plane, the variation of the function being determined
by the appropriate variations of the boundary of the
image of this domain.
The fundamental lemma of the method of boundary
variation. Let D be a domain in the w-plane and let
the complement!::. of D in the extended plane consist of
a number of continua. Let r be a continuum in !::. and
let there exist on r an analytic function s(w )=O such
that for any point Wo E r and for any function F( w)
that can be represented as
AJp2
F(w) = w+A o+ _ _ +O(p3) (*)
W-Wo
and that is univalent in D, the inequality
Re{AJs(wo)}+O(p);;;' 0
be valid, and suppose that the estimate of the residual
term in (*) is uniform in all closed subdomains of D. r
will then be an analytic curve that may parametrically
be represented by means of the function w=w(t) of the
real parameter t. This parameter may be so chosen that
r satisfies the differential equation
[c;; fS(W)+l = o.
This result makes the important role played by qua-
464
dratic differentials (cf. Quadratic differential) in the
solution of extremal problems in the theory of
univalent functions clear, since s(w) proves to be a
meromorphic function in many applications. In certain
cases it follows from the conditions of the problem that
the appropriate poles of s(w) belong to the boundary of
the extremal domain, and it is shown by the fundamen-
tal lemma of the method of boundary variation that the
boundary of this domain belongs to the union of the
closures of the critical trajectories of the quadratic dif-
ferential
In a number of extremal problems, the fundamental
lemma not only yields qualitative results, but also gives
sufficient information for the determination of the
boundary of the extremal domain, and hence for the
complete solution of the problem.
The following results were obtained by means of the
method of boundary variation: Qualitative results in
the coefficient problem for the class S; in the problem
of the maximum of the n-th diameter in a family of
continua of a given capacity; the solution of a number
of extremal problems of univalent conformal mappings
of doubly-connected domains; distortion theorems for
multiply-connected domains, which at the same prove
existence theorems of univalent conformal mappings of
a given multiply-connected domain onto canonical
domains, etc.
References
[I) SCHIFFER, M.: 'A method of variation within the family of sim-
ple functions', Proc. London Math. Soc., Ser. 244 (1938), 432-
449.
(2) ScHIFFER, M.: 'Some recent developments in the theory of con-
formal mappings', in R. Courant (ed.): Dirichlet's principle.
Conformal mapping and minimal surfaces, Interscience, 1950.
(3) ScHIFFER, M.: 'Applications of variational methods in the
theory of conformal mapping', in Calculus of variations and its
applications, Proc. Symp. Appl. Math., Vol. 8, Arner. Math.
Soc. & McGraw-Hill, 1958, pp. 93-113. GV K ' .
.. uz mma
Editorial comments. The 'fundamental lemma of the
method of boundary variation' is also known as Schiffer's
theorem.
References
[A 1] DUREN, P.L.: Univalent functions, Springer, 1983, Chapt.1 O.
AMS 1980 Subject Classification: 30C45, 30C70
BOUNDED OPERATOR - A mapping A of a topo-
logical vector space X into a topological vector space Y
such that A (M) is a bounded subset in Y for any
bounded subset M of X. Every operator A: X ~ Y, con-
tinuous on X, is a bounded operator. If A: X ~ Y is a
linear operator, then for A to be bounded it is suffi-
cient that there exists a neighbourhood U C X such that
A (U) is bounded in Y. Suppose that X and Yare
normed linear spaces and that the linear operator
 BRACHISTOCHRONE

A : X ~ Y is bounded. Then is finite.

y = sup II Ax II < 00.
2) A bounded set in a topological vector space E
Ilxll ;;;] (over a field k) is a set B which is absorbed by every
This number is called the norm of the operator A and is neighbourhood U of zero (i.e. there exists an aEk such
denoted by I A II. Then that B Call).
M.l. Voitsekhovskii
II Ax II ~ II A I . II x II, AMS 1980 Subject Classification: 46AXX, 54E35
and IAI is the smallest constant C such that
II Ax I ~ C IxI BOUNDEDLY-COMPACf SET in a topological linear
space X - A set M with the property that the closure N
for any x EX. Conversely, if this inequality is satisfied,
then A is bounded. For linear operators mapping a of every bounded subset N CM is compact and is con-
normed space X into a normed space Y, the concepts tained in M (for a normed space in the strong (resp.
of boundedness and continuity are equivalent. This is weak) topology this is equivalent to the compactness
(resp. weak compactness) of the intersections of M with
not the case for arbitrary topological vector spaces X
balls). A convex closed set in a normed space is
and Y, but if X is bomological and Y is a locally con-
boundedly compact if and only if it is locally compact.
vex space, then the boundedness of a linear operator
Boundedly-compact sets have applications in the theory
A : X ~ Y implies its continuity. If H is a Hilbert space
of approximation in Banach spaces; they have the pro-
and A: H ~H is a bounded symmetric operator, then
perty that an element of best approximation exists. A
the quadratic form <Ax, x> is bounded on the unit
barrelled topological linear space which is boundedly
ball K, ={x: I x 11':;;;1}. The numbers
compact (in itself) in the weak (resp. strong) topology is
f3 = sup <Ax, x> and a = Ilxllinf <Ax, x> a reflexive (resp. Montel) space. A normed space which
Ilxll;;;] <;;]

are called the upper and lower bounds of the operator A. is boundedly compact is finite-dimensional.
The points a and f3 belongs to the spectrum of A, and References
[1) KLEE, V.L.: 'Convex bodies and periodic homeomorphisms in
the whole spectrum lies in the interval [a, f3]. Examples
Hilbert space', Trans. Amer. Math. Soc. 74 (1953), 10-43.
of bounded operators are: the projection operator (pro- (2) EDWARDS, R.E.: Functional analySiS: theory and applications,
jector) onto a complemented subspace of a Banach Holt, Rinehardt, Winston, 1965.
L.P. Vlasov
space, and an isometric operator acting on a Hilbert
space. AMS 1980 Subject Classification: 46AXX
If the space X and Y have the structure of a partially
ordered set, for example are vector lattices (cf. Vector BOURGET FUNCTION - The function In,k(Z) which
may be defined as a generalization of the integral
lattice), then a concept of order-boundedness of an
representation of the Bessel functions
operator can be introduced, besides the topological
boundedness considered above. An operator A: X ~ Y
is called order-bounded if A (M) is an order-bounded set
In,k(Z) = 2~)t-n-] [t++r exp { ~Z[t-+]}dt,
in Y for any order-bounded set M in X. Examples: an where n is an integer and k is a positive integer. The
isotone operator, i.e. an operator such that x '0' implies integration contour makes one counter-clockwise tum
Ax':;;;Ay. around the coordinate origin. In other words,
References 1 '1T
=-

[1) LYUsTERNIK, L.A. and SOBOLEY, V.I.: E/emente der Funk- In,k(Z) j(2 cos 81 cos(n8-z sm8) d8,
tiona/analysis, Akademie Verlag, 1955 (translated from the 7r 0
Russian).
(2) RUDIN, W.: Functiona/ analysis, McGraw-Hill, 1973. In,O(z)=.Jn(z) is a cylinder function of the first kind. So
[3) BIRKHOFF, G.: Lattice theory, Arner. Math. Soc., 1967. named after 1. Bourget [1], who studied the function
v.I. Sobolev with a view to various applications in astronomy.
Editorial comments.
References
References [1) BoURGET, J.: 'Memoire sur les nombres de Cauchy et leur
[A1) TAYLOR, A.E. and LAY, D.C.: Introduction to functional application a divers problemes de mecanique celeste', J. Math.
analysis, Wiley, 1980. Pures App/. (2) 6 (1861), 32-54.
[A2) ZAANEN, A.c.: Riesz spaces, 2, North-Holland, 1983. (2) WATSON, G.N.: A treatise on the theory of Besse/functions,
Cambridge Univ. Press, 1952, Chapt. 10.
AMS 1980 Subject Classification: 47 A05, 47855 v.I. Pagurova
AMS 1980 Subject Classification: 33A40
BOUNDED SET - I) A bounded set in a metric
space X (with metric p) is a set A whose diameter BRACIDSTOCHRONE - A curve of fastest descent.
8(A) = sup p(x, y) The problem of finding such a curve, which was posed
x,yEA

465
BRACHISTOCHRONE

by G. Galilei [1], is as follows: Out of all planar curves

connecting two given points A and B and lying in a
common vertical plane (B below A) to find the one
along which a material point moving from A to B
solely under the influence of gravity, would arrive at B
within the shortest possible time. The problem can be
reduced to finding a function y (x) that constitutes a cr,
minimum of the functional

J(y) = f -if
u
1 +y'2 dx,
2gy
Fig. 1. Fig. 2.

If this is the identity permutation, w is called a coloured

where a and b are the abscissas of points A and B. The
brachistochrone is a cycloid with a horizontal base and
with its cusp at the point A.
References
[I] GALILEI, G.: Unterredungen und mathematische Demonstra-
tionen aber zwei neue Wissenszweigen, die Mechanik und die
Fallgesetze betreffend, W. Engelmann, Leipzig, 1891 (translated
from the Italian and the Greek).
[2] LAVRENT'EV, M.A. and LVUSTERNIK, L.A.: A course in varia-
tional calculus, Moscow-Leningrad, 1950 (in Russian).
L.P. Kuptsov
Editorial comments. The brachistochrone problem is
usually ascribed to Johann Bernoulli, cf. [A1], pp. 350 ff. A
treatment can be found in most textbooks on the calculus of
variations, cf. e.g. [A2].
References
[A1] ROUSE BALL, W.W.: A short account of the history of
mathematics, Dover, reprint, 1960.
[A2] WEINSTOCK, R.: Calculus of variations, Dover, reprint, 1974.
AMS 1980 Subject Classification: 49-XX, 51 M15
BRAID THEORY - The branch of topology and alge-
bra concerned with braids, the groups formed by their
equivalence classes and various generalizations of these
groups [1].
A braid on n strings is an object consisting of two
parallel planes Po and P I in three-dimensional space
R3 , containing two ordered sets of points
aJ, ... ,anEP O and b l , . . . ,bnEP I , and of n simple
non-intersecting arcs II, ... ,In' intersecting each paral-
lel plane Pt between Po and P I exactly once and join-
ing the points {a;} to {b;}, i = 1, ... ,n. It is assumed
that the a;'s lie on a straight line La in Po and the b;'s
on a straight line Lb in P I parallel to La; moreover, bi
lies beneath aj for each i (see Fig. 1). Braids can be
represented in the projection on the plane passing
through La and L b; this projection can be brought into
general position in such a way that there are only fin-
itely many double points, each two of which lie at dif-
ferent levels, and the intersections are transversal.
The string Ii of a braid w joins ai to bki and so
defines a permutation
S'" = [k\ ::: :n] .
(or pure) braid. The transposition (i i + 1) corresponds
to a simple braid (Ji (see Fig. 2).
On the set of all braids on n strings with fixed Po,
P J, {aj}, {b;}, one introduces the equivalence relation
defined by homeomorphisms h : II ~ II, where II is the
region between Po and PI, which reduce to the identity
on Po U PI; it may be assumed that h(Pt)=Pt . Braids
a and {3 are equivalent if there exists a homeomorphism
with the above properties such that h(a)={3.
The equivalence classes - which are still called
braids - form the braid group B(n) with respect to the
operation defined as follows. Place a copy II' of the
domain II above another copy II", in such a way that
P~ coincides with p'J, a;' with b;, and then compress
II' U II" to half its 'height'. The images of the braids
w' E II' and w" E II" produce a braid w' w", with. string Ii
obtained by extending f; with (~i' where k j E S'" (i). The
identity braid is the equivalence class containing the
braid with n parallel segments; the inverse w-I of a
braid w is defined by reflection in the plane P I / 2 For
the condition ww- I = see Fig. 3.
Fig. 3.
The mapping w~S'" defines an epimorphism of B(n)
onto the group S(n) of permutations of n elements, the
kernel of this epimorphism is the subgroup K(n) of all
pure braids, so that one has an exact sequence
1 ~ K(n) ~ B(n) ~ S(n) ~ 1.
The braid group B(n) has two principal interpreta-
tions. The first, as a configuration space, is obtained by
identifying the planes P t via vertical projection onto
Po, under which the images of the points ail =/i Pt, n
considered as t varies from 0 to 1, form the trace of an

466
 BRAID THEORY

isotopy $~ of the set U aj along Po; one has The braids OJ, 1~i ~n -1, are the generators of the
$r< U aj)= U aj. Consider the space of unordered group B(n), i.e. w=Ok, . . . 0k with m ,

sequences G(n) of n pairwise distinct points of the (Ji(Jj = (Jj(Ji if I i - j I > 1, }

plane; then each braid corresponds in one-to-one (Ji(Ji+1(Ji = (Ji+1(Ji(Ji-], l';;;i';;;n -2. (1)
fashion to a class of homotopy loops in this space, and
It turns out that (1) is a presentation of B(n) (see Fig.
one has an isomorphism
4).
{3: B(n) ~ 71"1 G(n).
There exists a splitting exact sequence (obtained from
For pure braids one has an analogously constructed the locally trivial fibration F(n )~F(n - 1) with fibre
isomorphism R2 \ (aI, ... ,an-I:
a: K(n) ~ 71"1 F(n),
1 ~Fn-1 ~K(n)~K(n-l)~ 1,
where F(n) is the space of ordered sequences of n dis-
tinct points of the plane, so that K(n) can be identified which leads to the normal series
with the subgroup corresponding to the covering K(n) = An :l ... :l A1 :l Ao = Fn- 1
p:F(n)~ G(n) = F(n)jS(n). with free fractors Aj / Aj _ I, where Aj has a 'com-
The second interpretation, as a homeotopy group, is ponent' Un - j isomorphic to K(n -i -1). Each element
obtained by extending the isotopy $~ to an isotopy $~ wEB(n) can be expressed uniquely in the form
of the plane Po that coinci~es with the identity outside
some disc, and such that $~ =id. For each t, two such where 71"", is a selected representative for S'" in B(n)
extensions differ by a homeomorphism which is the
identity at the points ail. A braid uniquely determines a
and Wi EAn - i + I n
Uj The reduction of a braid to this
form is known as its dressing. This solves the word
component of the space of homeomorphisms Y(n) of (identity) problem in B(n).
the plane which map the set U aj onto itself, and one A presentation of K(n) is as follows: generators (see
has an isomorphism Fig. 5):
'I: B(n) ~ 71"0 Y(n).

To each homeomorphism hEY corresponds an auto-

morphism of the free group of rank n : relations:
Fn = 71"1 (R2 \ U aj), defined up to an inner automor-
if r.;;;s<i<j
phism, which in turn yields a homomorphism
or i<r<s<j;
B(n)~OutFn =AutFn /InnFn The elements of the
image are called braid automorphisms of the free group. ArjAij = Arj if i =s; (2)

In particular, corresponding to the braid OJ one has an AijAsjAijAij = Asj if r=i<j<s;

automorphism
if }=I=i, i + 1 ({x;} is a set of generators of Fn). Any
braid automorphism a possesses the following
properties:
up to an inner automorphism (for the meaning of A j ,
see below); these properties characterize braid automor-
phisms.
This presentation may be obtained as a presentation of
the kernel of the natural homomorphism into S(n) of
the abstract group B(n) defined by the present(!tion (1)
~th the aid of the Schreier system rr;=2Mjkj,
J~k/~n, where Mji=oj_1 ... OJ.
The centre of B(n) is the infinite cyclic group gen-
erated by the element (01 . . . 0n)n. The commutator
group B'(n) coincides with B"(n) for n~5; B'(3) is
isomorphic to the free group of rank 2, and B' (4) is
isomorphic to semi-direct product of two such groups.
The quotient group modulo the commutator subgroup
is an infinite cyclic group, generated by the images of
OJ. There are no elements of finite order in B(n). The

group K(n) is mapped onto itself by endomorphisms

with non-Abelian image. In particular, K(n) B' (n) is n
a fully-characteristic subgroup of B(n), and also of
K(n) (see [15]).
Solving the conjugacy problem in B(n) is much more
Fig. 4. Fig. 5. complicated than solving the word problem. There is a

467
BRAID THEORY
unique Garside normal form of a braid, where
~=(Ol ... 0n-I)(OI ... 0n-2)' .. (102)01 is what is
known as a Garside element and 0 is a positive braid,
i.e. a braid the representation of which in terms of the
elements 0i has positive indices. With any braid '" one
can associate, using finitely many operations depending
on i (conjugation with certain elements, choice of ele-
ments of maximum degree, etc.), a certain set of words
~("'), from which one selects a word in normal form
~ + T with minimal T. This is a so-called upper form of
the braid "'. It turns out that two braids are conjugate
if and only if they have a same upper form (see [7]).
The Burau representation of the braid group B(n) in
the group of matrices over the ring of polynomials in
one variable with integer coefficients is defined by the
correspondence
Ii-I 0 0
0 i-I I
b(w): 0i ~
0 1 0 0
0 0 0 In-k-I
where h is the identity matrix of order k. The matrix
b(",)- In is the reduced Alexander matrix (see Alex-
ander invariants) of the link obtained by closing the
braid '" (see below). For a pure braid one obtains the
full Alexander matrix from the analogous Gassner
matrix. The problem of whether these representations
are faithful is still (1982) unsolved (see [2]).
The fact that the spaces F(n) and G(n) are aspherical
makes it possible to evaluate the homology of braid
groups.
The homology of K(n) (see [16]). Homologically,
K(n) coincides with the product of unions of circles in
which the number of circles runs from one through
n - 1. The homology ring is isomorphic to the exterior
graded ring generated by the one-dimensional elements
"'U="'Ji' l:s;;,i~j:s;;,n, with relations
WklWlm +WlmWmk +WmkWkl
As "'kl one can take the forms
1 dzk-dz i
2'TTi Zk -ZI '
corresponding to passage along the diagonals Zk =z,.
The homology of B(n) (see [8], [12]). The homomor-
phism B(n )~S(n) can be extended to an imbedding
Sen )~O(n); the induced homomorphism of cohomol-
ogy spaces H* (O(n ~H* (B(n is epimorphic, i.e. the
cohomology spaces mod 2 of the group B(n) generated
by the Stiefel- Whitney classes.
There is a natural mapping of G(n) into 0 2S2, the
two-fold loop space of S2, i.e. the space of spheroids
S2~S2 (choose small discs about n points, then map
these discs canonically with degree 1 into a sphere,
mapping the entire complement onto a point). This
468
",=~mo, mapping (see [14]) establishes a homology equivalence
of the limit space G(oo) and (02S2)o (the subscript
indicates that one chooses the component of spheroids
of degree 0). As to the unstable homology groups of
B(n), it has been proved [16] that they are finite, stabil-
ize at height n and satisfy the recurrence relation
H i B(2n+I)=H i B(2n). There is a description of the
calculation of these groups [17].
Applications and generalizations. I) A closed braid is a
link (an n-component knot) in R3 each component of
which transversally cuts out half-planes bounded by the
same straight line: the axis I of the closed braid (see
Fig. 6).
1
~
0
0
Fig. 6.
Fig. 7.
w
A braid '" generates a closed braid (the closure of w)
in the following way. Consider a cylinder with bases on
Po and PI, the interior of which contains w. Let this
cylinder be deformed in R3 so that its elements become
circles with centres on I, its bases coincide and each
w
point ai coincides with bi' Then is the union of the
strings Ii. Conversely, every link in R3 can be
represented by a closed braid. To equivalent braids
= O. correspond isotopic links and, moreover, conjugate
braids yield isotopic links. The converse is false, since a
link may be represented by braids with different
numbers of strings. In addition, the braids Wan -I and
WO;~I are not conjugate in B(n), but they correspond
to isotopic links. If two closed braids are equivalent as
links, they can be derived from one another by a chain
of elementary transformations of two types (see Fig. 7).
These operations are interpreted in terms of presenta-
tions of the link group, thus yielding a reformulation of
the isotopy problem for links as a question concerning
the system of groups B(n). A presentation of the link
w
group of has the form
VI, ... ,Yn: Yi =AiYkjAi- I },
where the relations are defined by braid automorphisms
b W Conversely, any such relation defines a braid.

2) If one cuts a surface of genus g with g non-
intersecting cuts so as to obtain a sphere with 2g holes,
then the homeomorphisms of this sphere with holes
that fix points on the edges of the holes define
homeomorphisms of the surface which fix the cuts and
are themselves defined up to isotopy by the elements of
the group K(2g). This yields a representation of the
braid group in the homotopy group of the surface.
Similarly one constructs a representation of B(2g).
These representations are used in studying Heegaard
diagrams of three-dimensional manifolds (cf. Three-
dimensional manifold).
3) By identifying R2 with the complex plane CI and
associating with any unordered set of n points in the
plane a polynomial of degree n having these points as
roots, one can identify G(n) with the space of polyno-
mials with non-zero discriminants. This has made it
possible to obtain several results concerning the non-
representability of algebraic functions by superpositions
of functions in fewer variables (see [16]).
4) Configuration spaces for arbitrary spaces X are
defined in analogy with G(n) and F(n), with R2
replaced by X. The fundamental groups of these spaces,
B(X) and K(X), are called the braid group of the space
X and the pure braid group of the space X, respectively.
For a manifold M n of dimension exceeding 2,
17IFn(X);::::::II~=117,M(i), and this group is of no
interest. For two-dimensional manifolds, one has a
natural imbedding of B(n) and K(n) into Bn(M2) and
Kn(M2) induced by an imbedding R2 CM2. If M2 is
neither a sphere nor a projective space, one obtains an
exact sequence
e n
1~ 7TI K(2) ~ 7TI Kn(M2) ~ II7T;Mfi) ;
;=1
for the sphere, the homomorphism e is an epimor-
phism, obtained by adding to (l) the single relation
5) If p: X ~ Y is a k-sheeted covering, then p -I ex,
where ex is a loop in Y, is a loop in the configuration
space X, and this defines a homomorphism
171 Y ~Bk(X) which strengthens the monodromy of the
covering and has applications in algebraic geometry.
6) Let VC be the complexification of a real vector
space V and let W be a finite irreducible group gen-
erated by reflections acting in V (hence also in VC).
Let s, be generating reflections in the planes P, C V and
let D be their union. Finally, let VC / D = Yw and let
X w be the quotient space. The groups 171 Y w and 171 X w
are called Brieskorn groups and constitute natural gen-
eralizations of K(n) and B(n). If ord(s,sj) = m'j' then
17 I X W has a presentation of the form
BRANCH INDEX
where the number of factors on each side is equal to
mij (0, here corresponds to a Weyl chamber). It has
been proved for these groups that X w and Y w are
spaces of type K(17, I), and the conjugacy problem has
been solved. The spaces X w appear in algebraic
geometry as complements to the discriminant of versal
deformations of rational singularities (see [12], [13]).
References
[1] ARTIN, E.: 'Theory of braids', Ann. of Math. 48 (1947), 643-
649.
[2] BIRMAN, l.S.: Braids, links and mapping class groups, Princeton
Univ. Press, 1974.
[3] BURAU, W.: 'Ueber Zopfinvarianten', Abh. Math. Sem Univ.
Hamburg 9 (1932), 117-124.
[4] MARKov, AA: 'Foundations of the algebraic theory of
braids', Trudy Mat. Inst. Steklov. 16 (1945) (in Russian).
[5] GASSNER, B.: 'On braid groups', Abh. Math. Sem Univ. Ham-
burg 25 (1961), 10-22.
[6] FADELL, E. and NEUWIRTH, L.: 'Configuration spaces', Math.
Scand 10 (1962), 111-118.
[7] GARSIDE, F.A: 'The braid group and other groups', Quart. J.
Math. 20, no. 4 (1969), 235-254.
[8] FUKS, D.B.: 'Cohomology of braid groups mod 2', Funktional.
Anal. i Prilozhen. 4, no. 2 (1970), 62-73 (in Russian).
[9] ARNOL'D, V.I.: 'On cohomology classes of algebraic functions
that are preserved under Tschimhausen transformations',
Funktional. Anal. i Prilozhen. 1 (1970), 84-85 (in Russian).
[10] GoRIN, E.A and LIN, V.YA.: 'Algebraic equations with con-
tinuous coefficients and some problems in the theory of
braids', Mat. Sb. 78 (1969),579-610 (in Russian).
[11] ARNOL'D, V.I.: 'On certain topological invariants of algebraic
functions', Trudy Moskov. Mat. Obshch. 21 (1970), 27-46 (in
Russian).
[12] BRIESKORN, E.: Matematika 18, no. 3 (1974), 46-59.
[13] BRIESKORN, E. and SAITO, K.: 'Artin Gruppen und Coxeter
Gruppen', Invent. Math. 17 (1972), 245-271.
[14] DELIGNE, P.: 'Les immeubles des groupes de tresses
generalises', Invent. Math. 17, no. 4 (1972), 273-302.
(15] LIN, V.YA.: 'Representations of permutation braids', Uspekhi
Mat. Nauk. 29, no. 1 (1974), 173-174 (in Russian).
(16] ARNOL'D, V.I.: 'The ring of cohomology groups of crossed
braids', Mat. Zametki 5, no. 2 (1969), 227-231 (in Russian).
[17] LIN, V.YA.: 'Artin braids and related groups and spaces', Itogi
Nauk. i Tekhn. Algebra Topol. Geom. 17 (1979), 159-227 (in
Russian). A .V. Ch ernavsk"II
AMS 1980 Subject Classification: 57M25
BRANCH INDEX - The sum V = ~(k - 1) of the
orders of the branch points (cf. Branch point) of a com-
pact Riemann surface S, regarded as an n-sheeted cov-
ering surface over the Riemann sphere, extended over
all finite and infinitely-distant branch points of S. The
branch index is connected with the genus g and
number of sheets n of S by:
V = 2(n+g-l).
See also Riemann surface.
References
[1] SPRINGER, G.: Introduction to Riemann surfaces, Addison-
Wesley, 1957.
E.D. Solomentsev
Editorial comments.
469
BRANCH INDEX

References Il(z] ; r) is again obtained, this is also true of all ele-

[A1] MUMFORD, D.: Algebraic geometry I. Complex projective ments of the branch (cf. Branch of an analytic function)
varieties, Springer, 1976.
of f (z) defined in V by the element Il(z] ; r). In such a
AMS 1980 Subject Classification: 14XX, 30F10 case a is a branch point of finite order k - I of this
branch. In a punctured neighbourhood V of a branch
BRANCH OF AN ANALYfIC FUNCflON - The point a of finite order this branch is represented by a
result of analytic continuation of a given element of an generalized Laurent series, or Puiseux series:
analytic function represented by a power series
00 (I)
n(a; r) = ~cv(z -at v= -00
v=O

with centre a and radius of convergence r >0 along all If a = 00 is an improper branch point of a finite
possible paths belonging to a given domain D of the order, then the branch of f (z) is representable in some
complex plane C, a ED. Thus, a branch of an analytic neighbourhood V' = {z: 1 z 1 > p} by an analogue of
function is defined by the element Il(a; r) and by the the series (I):
domain D. In calculations one usually employs only
(2)
single-valued, or regular, branches of analytic functions, v= -00

which exist not for all domains D belonging to the

domain of existence of the complete analytic function. The behaviour of the Riemann surface R of f (z) over
For instance, III the cut complex plane a branch point of finite order a is characterized by the
D =C\ {z =x: - oo";;;x,,;;;o} the multi-valued analytic fact that k sheets of the branch of f (z) defined by the
function w = Ln z has the regular branch element Il(z]; r) come together over a. At the same
time the behaviour of other branches of R over a may
w = Lnz = In I z I +iargz, I argz 1<'11', be altogether different.
which is the principal value of the logarithm, whereas If the series (l) or (2) contains only a finite number
in the annulus D = {z: 1< 1 z 1 <2} it is impossible to of non-zero coefficients bv with negative indices v, a is
isolate a regular branch of the analytic function an algebraic branch point or an algebraic singular point.
w=Lnz. Such a branch point of finite order is also characterized
References by the fact that as z~a in whatever manner, the values
[I] HURWITZ, A. and COURANT, R.: Vorlesungen iiber allgemeine of all elements of the branch defined by Il(z]; r) in V
Funktionentheorie und elliptische Funktionen, I, Chapt. 3; or V' tend to a definite finite or infinite limit.
2, Chapt. 4, Springer, 1964.
[2] MARKUSHEVICH, A.I.: Theory offunctions of a complex variable,
Example: f(z) = z ]1 k , where k> I is a natural
2, Chelsea, 1977 (translated from the Russian). number, a =0,00.
E.D. Solomentsev If the series (l) or (2) contain an infinite number of
AMS 1980 Subject Classification: 30A05, 30B40 non-zero coefficients bv with negative indices v, the
branch points of finite order a belong the class of
BRANCH POINT, singular point of multi-valued char- transcendental branch points.
acter - An isolated singular point a of an analytic func- Example: f(z)=exp(l/z)]lk, where k>1 is a
tion f (z) of one complex variable z such that the ana- natural number, a =0.
lytic continuation of an arbitrary function element of Finally, if it is impossible to return to the initial ele-
f (z) along a closed path which encircles a yields new ment after a finite number of turns, a is said to be a
elements of f (z). More exactly, a is said to be a branch logarithmic branch point or a branch point of infinite
point if there exist: I) an annulus order, and is also a transcendental branch point.
V={z: 0<1 z-a I<p} in whichf(z) can be analyti- Example: f(z)=Lnz, a =0,00.
cally extended along any path; 2) a point z] E V and Infinitely many sheets of the branch of f (z) defined
some function element of f (z) represented by a power by the element Il(z]; r) come together over a loga-
senes rithmic branch point.
x
n(zl;r) = ~cvCz-zd\ In the case of an analytic function of several complex
v =-=0
variablesf(z), z=(z], ... ,zn), n;;.2, a point a of the
with centre z] and radius of convergence r>O, the ana- space en or cpn is said to be a branch point of order m,
lytic continuation of which along the circle l";;;m";;; 00, if it is a branch point of order m of the,
1z -a 1= 1z] -a I, going around the path once in, generally many-sheeted, domain of holomorphy of f (z).
say, the positive direction, yields a new element Unlike in the case n = I, branch points, just like other
rI'(z] ; r') different from Il(z]; r). If, after a minimum singular points of analytic functions (d. Singular point),
number k > I of such rounds the initial element cannot be isolated if n ~2.

470

References
[I] MARKUSHEVICH, A.I.: Theory offunctions of a complex variable,
2, Chelsea, Chapt. 8 (translated from the Russian).
[2] FUKS, B.A.: Theory of analytic functions of several complex vari
ables, Amer. Math. Soc., 1963 (translated from the Russian).
E.D. Solomentsev
AMS 1980 Subject Classification: 30A05, 30B40
BRANCIllNG OF SOLUTIONS, bifurcation of solu-
tions, of non-linear equations - A phenomenon that
causes a given solution of a non-linear equation to
disappear completely or to become a number of solu-
tions if minor changes in the parameters of the equa-
tion are introduced. More precisely, let the non-linear
equation
F(x, A) = (*)
with parameter A (which need not be numerical) have
the solution Xo for a given value Ao of the parameter.
Then, if the values of A are close to Ao, equation (*)
may have more than one solution X(A) close to xo. One
then speaks of branching (bifurcation) of the solution xo,
and the pair (xo, Ao) is said to be a branching (bifurca-
tion) pOint of equation (*).
Example: The equation x 2 -A=O, where x and A are
complex variables, has the branching point
(xo, Ao)=(O, 0) because there exists a two-valued solu-
tion x = -..[i;, i.e. if A*O is small, the solution x =0 (for
A=O) branches out into two small non-trivial solutions.
The modem theory of branching of solutions is
based on the ideas of A.M. LY3punov [1] and E.
Schmidt [2] and has mostly been developed for non-
linear equations in Banach spaces.
Let Eland E2 be complex Banach spaces, x EE I, let
A be a complex variable, and let F(x, A) be a non-linear
operator which, together with its Frechet derivative
FxCx, A), is continuous in a neighbourhood Q of the
point (xo, Ao). Let F(x, A) map Q into a neighbourhood
of zero of the space E2 such that F(xo, Ao)=O, and
suppose that FxCxo, Ao) B is a Fredholm operator.
It is required to find, in the sphere II x - Xo " <r,
where the radius r is sufficiently small, all the solutions
of equation (*) that are continuous if I A- Ao I <p,
where p is also sufficiently small. In other words, this is
a problem of local extension of the solution Xo in the
parameter A. If the inverse operator B- 1 exists, the
problem has a unique solution X(A), and X(Ao)=Xo. If,
on the other hand, B- 1 does not exist, then the null
space N (B) of B has dimension n;;;;'1. In such a case
the problem may be reduced to an analogous finite-
dimensional problem. Let P denote the projection (d.
Projector) of E I on N (B) and let / - Q denote the pro-
jection of E2 onto the domain of values of B, where /
is the identity operator. Equation (*) may be written as
the system
BRANCHING OF SOLUTIONS
(I-Q)F(xo+u+v, A) = 0, }
QF(xo +u +v, A) = 0,
where u =(/ - P)(x - xo), v =P(x - xo). The implicit
operator u = f (v, A) is found from the first equation of
the system. Its substitution into the second equation of
the system yields the equation
QF(xo + j(v, A)+V, A) = 0,
from which v may be found; it is said to be the branch-
ing (bifurcation) equation. The complete solution of the
problem - to find in a sphere" v "<r of sufficiently
small radius r all solutions v (A) of the branching equa-
tion that are continuous for I A-Ao I <p, where p is
sufficiently small - yields the complete solution of the
initial problem, since all its solutions can be
represented in the form
X(A) = Xo +V(A)+JIV(A), Aj,
where v is a certain solution of the branching equation.
Let F(x, A) be an analytic operator in Q. The choice
of bases in the n-dimensional subspaces PEl =N(B)
and QE 2 makes it possible to write the branching
equation in the form of the system
2"j(~!, ... '~n,A) = 0, i=l, ... ,n,
where 2;, i = 1, ... ,n, is an analytic function at the
point (0, ... ,0, Ao), while all partial derivatives
a.!l'; / a~j vanish at this point. This system can be stu-
died using the theory of exclusion, the method of the
Newton diagram and other methods [3], [4], [5]. If
n = 1, the complete analysis is effected by the method
of the Newton diagram. As regards the study of the
branching equation, i.e. of the original problem as well,
the following three cases are the only ones possible: a)
the problem has no solution; b) the problem has a fin-
ite number of solutions, all of which can be represented
by converging series of integer or rational powers of the
difference A- Ao; or c) the problem has a finite number
of families of solutions, each one of which depends on
a finite number of small free parameters and, possibly,
on the finite number of solutions given in b).
For case b) to occur, it is sufficient for Xo to be an
isolated solution of the equation F(x, Ao)=O. In case
b) it is expedient to look for solutions by the method of
indefinite coefficients in the form
00
X(A) = xo+ ~Xk(A-Aot/P,
k=!
where Xk are the coefficients to be found, while the
possible values of p may be preliminarily found using
the branching equation. The substitution of such a
series in (*) yields a recurrent system for finding
x I , X 2, . . .. Problems of the type
BXk = H (x I, ... ,Xk -I) are obtained, and each Xk is
determined up to n arbitrary constants, which follow
471
BRANCHING OF SOLUTIONS
from the stipulated solvability of the successive equa-
tions. All resulting series converge in some neighbour-
hood of Ao. An estimate from below of the radius of
the neighbourhood may be obtained by the construc-
tion of majorants [6].
For case c) to occur, Xo must be a non-isolated solu-
tion of the equation F(x, Ao)=O. Here the application
of the method of indefinite coefficients may result in
divergent series (formal solutions). If the problem is
invariant with respect to a continuous group of linear
operators in E I, the use of group considerations will
make it possible in several cases to reduce the number
of the equ?tions and of the unknowns involved, and
thus simplify the problem or even reduce it to case b)
[7], [8].
Equation (*) may also have solutions that are defined
for A=Ao only. Such solutions are possible only if Xo is
a non-isolated solution of the equation F(x, Ao)=O;
they are found using the branching equation for A=Ao.
The determination of all its multi-parameter families of
solutions yields all the solutions of equation (*) with
A=Ao.
If the spaces E I and E 2 are real, the branching
equation is studied in the complex domain, after which
real solutions are taken. Some of them may prove to be
defined in half-neighbourhoods of Ao.
This procedure is also partly applicable if F(x, A) is
a sufficiently smooth operator, B is a Noetherian opera-
tor and the parameter A is an element of another
Banach space E (branching points in E may be
replaced by lines and surfaces). It is also used in the
study of certain related problems: the problem of find-
ing large solutions (equation (*) can have solutions
X(A)~OO as A~Ao), the problem of branching of eigen
values and eigen elements of linear operators, etc. [3].
The special case when
E 1 = E 2 , F(x,;\) - x -~(x, ;\), ~(O,;\) - 0,
was also studied by topological and variational
methods and methods involving the use of cones in a
Banach space. The concept of a bifurcation point is
very important in this circle of problems. Problems
regarding branching of solutions which do not belong
to this scheme are also encountered. They include, for
example, degenerate differential equations [9], [10] and
problems involving long and solitary waves [11].
References
[I] LYAPUNOV, A.M.: 'On equilibrium figures, etc.', in Complete
works, Vol. 4, Moscow, 1959 (in Russian).
[2] ScHMIDT, E.: 'Zur Theorie der linearen und nichtlinearen
Integralgleichungen III', Math. Ann. 65 (1908),370-399.
[3] VAlNBERG, M.M. and TRENOGIN, VA: Theory of branching of
solutions of non-linear equations, Noordhoff, 1974 (translated
from the Russian).
[4] VAlNBERG, M.M. and TRENOGIN, VA: 'The methods of
Lyapunov and Schmidt in the theory of non-linear equations
472
and their further development', Russian Math. Surveys 17, no.
2 (1962), 1-60. (Uspekhi Mat. Nauk 17, no. 2 (1962), 13-75)
[5] KRAsNOSEL'sKIi, M.A., ET AL.: Approximate solution of operator
equations, Moscow, 1969 (in Russian).
[6] AKHMEDOV, K.T.: 'The analytic method of
Nekrasov- Nazarov in non-linear analysis', Uspekhi Mat.
Nauk 12, no. 4 (1957), 135-153 (in Russian).
[7] YUDOVICH, V.I.: 'Free convection and bifurcation', J. Appl.
Math. Mech. 31 (1967), 103-114. (Priklad. Mat. Mekh. 31
(1967), 101-111)
[8] LOGINov, B.V. and TRENOGIN, VA: 'On the application of
continuous groups in the theory of branching', Soviet Math.
Dokl. 12, no. 2 (1971), 404-408. (Dokl. Akad Nauk SSSR 197,
no. 1 (1971
[9] AKHMEDOv, K.T.: 'The Cauchy problem for a class of non-
linear equations in function spaces', Dokl. Akad. Nauk SSSR
115, no. 1 (1957), 9-12 (in Russian).
[IO] SIDOROV, N.A.: 'Branching of solutions of the Cauchy problem
for a certain class of nonlinear integra-differential equations',
Differential equations N. Y. 3, no. 9 (1967), 830-834.
(Differentsial'nye Uravneniya 3, no. 9 (1967
[11] 1'ER-K1uKoROV, A.M. and TRENOGIN, V.A.: 'Long wave solu-
tions for quasi-linear elliptic equations in an unbounded strip',
Differential Equations N. Y. 3, no. 3 (1967), 496-508.
(Differentsial'nye Uravneniya 3, no. 3 (1967 VA 'T'
. . 1.renogm
Editorial comments. For a classification of types of
bifurcation (branching) pOints see Singularities of differen-
tiable mappings. The technique of reducing solving equa-
tion (*) to a finite-dimensional problem, the bifurcation
equation, is often called the Lyapunov-Schmidt method.
References
[A1] CHOW, S.-N. and HALE, J.K.: Methods of bifurcation theory,
Springer, 1982.
AMS 1980 Subject Classification: 47H15, 58E07,
58F14
BRANCHING POINT (OF A MINIMAL SURFACE)
- A singular point of a minimal surface in which the
first fundamental form of the surface vanishes; this
means, in fact, that such a branching point can exist on
a generalized minimal surface only. This singular point
owes its name to the fact that in a neighbourhood of it
the structure of the generalized minimal surface resem-
bles that of the Riemann surface of the function w =zn,
n ~2, over the point z =0, i.e. there the generalized
minimal surface has a many-sheeted orthogonal projec-
tion onto some plane domain, in which the projection
of the branching point itself is an interior point with a
unique pre-image. In a neighbourhood of a branching
point (u =0, v =0) the coordinates (x,y, z) of the
minimal surface can be represented in the form
x+ry = awm+O(1 w Im+l),
z = Re(bwm+n)+O(1 w Im+n+l), w = u+iv,
where a = c(l + i)==O and b==O are two complex con-
stants, 1m c = 0; m ~ 2, and n ~ 1 are integers named,
respectively, the order and the index of the singular
point and u and v are intrinsic isothermal coordinates.
The following theorem was deduced on the basis of

this representation: If the numbers m + n and m are
coprime, then the minimal surface has (m -I)(m +n)
different lines of self-intersection issuing from the
singular point in different directions. There is a relation
between the genus of a complete minimal surface, the
number of its branching points, and the index of its
Gaussian mapping [1].
One distinguishes between two kinds of branching
points: false branching points and true (non-false)
branching points. False branching points are singulari-
ties of the mapping which defines the surface that can
be got rid of by re-parametrization (e.g. if r=r(w) is a
regular minimal surface, then the generalized minimal
surface r=r(w2) has a false branching point at w =0).
A true branching point represents a real singularity of
the surface itself, and has the following important pro-
perty: In a neighbourhood of a true branching point
the surface can be altered so that the new surface which
coincides with the original one outside the deformed
neighbourhood will have a smaller area than the origi-
nal surface (this holds for surfaces in R3 but is not true
in more-general settings: e.g. for the area-minimizing
surfaces in R4 ). The theory of generalized minimal sur-
faces with a branching point served as a base of the
general theory of branched immersions, developed for a
broad class of two-dimensional surfaces in R", n ~ 3,
[2].
References
[1) CHEN, Y.W.: 'Branch points, poles and planar points of
minimal surfaces in R3', Ann. of Matk 49, no. 4 (1948), 790-
806.
[2) NITZSCHE, J.C.C.: Vorlesungen iiber Minimalj1iichen, Springer,
1973.
I.Kh. Sabitov
AMS 1980 Subject Classification: 53A 10
BRANCIUNG PROCESS - A stochastic process
describing a wide circle of phenomena connected with
the reproduction and transformation of given objects
(e.g. of particles in physics, of molecules in chemistry,
of some particular population in biology, etc.). The
class of branching processes is singled-out by the fun-
damental assumption that the reproductions of indivi-
dual particles are mutually independent.
A time-homogeneous branching process JL(t) with one
type of particles is defined as a Markov process with a
countable number of states 0, I, ... , the transition
probabilities Pit) of which satisfy the additional
branching condition:
Pit) = ~ PIj,(t) .. P1},(t). (1)
}'+"'+},=}
The states 0, I, ... , of a branching process are inter-
preted as numbers of particles. The probability PiJ(t) is
equal to the probability
P{p.(t+to)=j Ip.(to)=i}
BRANCHING PROCESS
of transition from state i to state j during a time inter-
val of length t. The principal analytical tools of branch-
ing processes are the generating functions (cf. Generat-
ing function)
00
F(t;s) = ~P{p.(t)=nlp.(O)=1}sn. (2)
n=O
The equality
F(t +".; s) = F(t; F(".; s (3)
follows from the branching condition. In branching
processes with discrete time, t assumes non-negative
integral values, and it follows from (3) that F(t; s) is a
t-fold iteration of the function F(s)=F(I;s). Such a
process is sometimes called a Galton - Watson process.
In a continuous-time branching process it is assumed
that t E [0, (0) and that the right-side derivative
aF(t; s)
at
I = /(s)
1=0
exists. It follows from (3) that F(t; s) satisfies the dif-
ferential equation
aF~/S) =/(F(t;s (4)
and the initial condition F(O;s)=s.
If A =F' (I) and a =/(1) are finite, the mathematical
expectation EJL(t) of the number of particles JL(t) (under
the condition JL(O) = I) is A I for a discrete-time branch-
ing process and eat for a continuous-time branching
process. Depending on the values of the parameters A
or a, branching processes are subdivided into subcritical
(A < 1, a<O), critical (A = I, a =0) and supercritical
(A> I, a>O). The main characteristic governed by this
classification is the behaviour of EJL(t) as t~oo.
In what follows, the trivial cases F(s)=s and
f(s) 0, when
P{p.(t)=llp.(O)=l} - 1,
will not be discussed.
The probability of extinction is one in subcritical and
critical branching processes and is less than one in
supercritical processes. If EJL(t}ln JL(t)< 00, then, in the
subcritica1 range, the probability P{JL(tO} of survival
of the process behaves asymptotically, as t~oo, as
K EJL(t), where K is a positive constant. In the case of
critical branching processes with finite Ep.2(t), the
asymptotic behaviour as t~oo is
2
P{p.(tO} ~ Op.(t) ,
where DJL(t)=B t in discrete-time branching process
and DJL(t)=bt in continuous-time branching process
and where B = F" (I), b =f" (1). A more detailed study
of the asymptotic behaviour of the distribution of JL(t)
as t~oo shows that the conditional distribution law
473
BRANCHING PROCESS
S/(x) = p{ E{JL(t)11to} ~x lJL(tO}
converges weakly as t~oo towards a limit distribution
Sex) if certain moments of p.(t) are finite. In subcritical
branching processes the limit law S (x) is discrete, while
in other cases it is absolutely continuous. Of special
interest is the case of critical branching processes for
which the limit law S (x) is exponential:
Sex) = l-e- x , x;;;'O.
The distribution (6) is also a limit distribution for
near-critical branching processes. More exactly, if one
considers the class of generating functions F (s) or f (s)
with a bounded third derivative F'" (1), j''' (1) and
F"(lr~Bo>O,j"(1r~bo>O, then
lim S/(x) = l-e- x , x;;;.O,
/-->00
A-->I
where St(x) is defined by formula (5). The phenomena
occurring in near-critical processes at t~oo are said to
be transient.
Another model of branching processes is the
Bellman - Harris process, in which each particle has a
random lifetime with distribution function G(t). At the
end of its lifetime, the particle leaves behind n daughter
particles with probability qn;
00
~qn = 1.
n=O
The lifetimes and the number of daughter particles pro-
duced by individual particles are independent. Consider
a particle of age zero at the initial moment of time
t =0; then the generating function F(t;s) of the
number of particles p.(t) at the moment of time t, given
by formula (2), satifies the non-linear integral equation
F(t;s) = jh(F(t-u;sdG(u)+s(l-G(t,
o number of particles present in subsets of X, the branch-
where 00
h(s) = ~q.sn.
n=O
If G(t) is a degenerate distribution function, the
Bellman - Harris process is a discrete-time branching
process; if G(t) is an exponential distribution function,
one obtains a continuous-time branching process. In
the general case, the Bellman - Harris process is a non-
Markov branching process.
A branching process may also be complicated by the
dependence of the particles on their location in space.
For instance, let each particle execute an independent
Brownian motion in an r-dimensional domain G with
an absorbing boundary aGo A particle located inside
the domain G will be transformed within time f1t~O
with probability
Pn!:.t +o(!:.t), n#l,
474
(5)
into n particles, and each such particle begins to exe-
cute an independent motion along Brownian trajec-
tories starting from its point of genesis. Let /Lx,t(A) be
the number of particles in a set A at the moment of
time t, that number having been produced by one parti-
cle located at the point x E G at the initial moment of
time 0. The generating function
R(t, x, s( = Eexp [[InS(Y)J.!X.1 (dY )]
(6)
satisfies the quasi-linear parabolic equation
t:.H + feR) =0 (8)
with the initial condition H(O, x, s(=s(x) and the
boundary condition H (t, x, s( Ix->aG =0, where
axt
f1 = ~; = 1 a2 / is the Laplace operator, and
f(s)=~"Pnsn'Pl = - ~n*lPn.
It is assumed in general branching processes that the
propagating particles can be characterized by certain
parameters, which may be interpreted as age, location
of the particle in space, type, size or energy of the par-
ticle, etc. Such processes are studied with the aid of
generating functions or functionals, for which non-
linear differential and integral equations, generalizing
equations (4), (7) and (8), are deduced. One may give
the following general descriptions of such models of
branching processes. Let particles move, independently
of each other, in accordance with Markov's law in some
phase space X. It is assumed that the random lifetime
of a particle is a Markov time (d. Markov moment)
which depends on its trajectory. At the end of its life-
time the particle generates new particles, which become
distributed over the phase space X in accordance with
some probability law. The new particles evolve in a
similar manner, independently of each other. In the
(7) space of integer-valued measures determined by the
ing process is Markovian. However, branching
processes are often studied in simpler reduced spaces as
well. In such a case many of them become non-
Markovian.
In most of the models discussed above the subdivi-
sion of processes into subcritical, critical and supercriti-
cal processes retains its meaning. Many properties of
simple branching processes described by equation (4)
are also displayed by more complex systems. In partic-
ular, the limiting distribution of (5), for critical
processes, usually turns out to be an exponential distri-
bution (6) (in an appropriate norm).
Branching processes are used in calculations of vari-
ous real processes - biological, genetic, physical,
chemical or technological. In real processes the condi-
tion of independent motion of each particle is often not
met, and the generated daughter particles may interact
 RRA'.;CHI'\(i PROCESS. A(i F-I)"P"'.; I)"'.; I

with each other. This is true of many biological repro-
duction processes, in processes of propagation of epi-
demics (cf. Epidemic process), in bimolecular chemical
reactions, etc. However, the initial stages of such
processes may be calculated with the aid of suitably
chosen models of branching processes. This is done in
situations in which the number of active particles in the
space is relatively small; at such low concentrations
their mutual collisions are practically non-existent, and
the state of the system evolves by the collisions between
active particles and the particles of the medium. In epi-
demic processes, for example, sick individuals may be
At the end of its life the particle is transformed into k
daughter particles of age zero with a probability Pk(u)
if the transformation took place when the age attained
by the original particle was u. Let JL(t) be the number
of particles at the moment of time t. The generating
function F(t;x) of the probability distribution of JL(t)
for a process beginning with one particle of age zero
satisfies the equation
F(t;x) = jh(u, F(t-u;xdG(u)+x(l-G(t, (*)
o
where 00

considered as such 'active particles'. Phenomena related

h(u, x) = ~Pk(U)Xk.
k =0
to genetic mutations, for example, may be calculated Put
with the aid of branching processes. The branching pro-
cess with a finite number of particle types may serve as
a model in computing chain reactions; the branching 00 00

process with diffusion may serve as a model of neutron A = ja(u)dG(u), B = jb(u)dG(u).

processes in nuclear reactors. Phenomena occurring in o 0

showers of cosmic particles may also be studied with
the aid of branching processes. In telephony, the com-
putation of certain expectation systems may be reduced
to branching process models.
See also Branching process with a random medium;
Branching process with immigration; Branching
processes, regularity of.
An age-dependent branching process is said to be sub-
critical, critical or supercritical if A > 1, A = 1 and B >0,
or A> 1, respectively. The behaviour of the process as
t ~ 00substantially depends on its criticality. Subcritical
and critical processes die out with probability one, i.e.
lim P{JJ.(t)=O} = 1.
t~oo

The following results have been obtained for these

References
[I] SEWASTJANOW, B.A. [B.A. SEVAST'YANOV]: Verzweigungspro-
zesse, Akad. Wissenschaft. DDR, 1974.
[2] ATHREYA, K.B. and NEY, P.: Branchingprocesses, Springer,
1972.
B.A. Sevast'yanov
Editorial comments. Applications of branching process
in biology may be found in [A2]. The survey article [A3]
gives an up-to-date account of stable population theory and
balanced exponential growth, items that are also relevant in
biological applications.
processes [1]: asymptotic formulas for the moments
JL(t), necessary and sufficient conditions of extinction,
conditions of existence and uniqueness of a solution of
equation (*) and asymptotic formulas as t ~ 00 for
Q(t) = P{JJ.(tO}.
The limit distributions have also been determined. In
the critical case, as t ~ 00 :
00

j ua(u)dG(u)
References
[A1] ASMUSSEN, S. and HERING, H.: Branching processes, Q(t):::::: 2-"0_ _ __
Bt
Birkhauser, 1983.
[A2] JAGERS, P.: Branching processes with biological applications,
Wiley, 1975.
[A3] JAGERS, P. and NERMAN, 0.: 'The growth and composition of
p{ E(JJ.(t)12tO) <xlJJ.(tO}~ l-e- x , x>O.

branching populations', Adv. Appl. Probab. 16 (1984),221-

259.
AMS 1980 Subject Classification: 60GXX, 60J80,
60J85
BRANClDNG PROCESS, AGE-DEPENDENT - A
model of a branching process in which the lifetime of a
particle is an arbitrary non-negative random variable,
while the number of daughter particles depends on its
age at the moment of transformation. In the single-type
particle model each particle has a random duration of
life T with distribution function
P{T";;;;t} = G(t).
If h (u, x) is independent of u, the age-dependent
branching process is a ReUman - Harris process. The
model just described has been generalized to include
processes with several types of particles, and also to
processes for which a particle may generate new parti-
cles several times during its lifetime [2], [3].
References
[I] SEWASTJANOW, B.A. [B.A. SEVAST'YANOV]: Verzweigungspro-
zesse, Akad. Wissenschaft. DDR, 1974.
[2] SEVAST'YANOV, B.A.: 'Age-dependent branching processes',
Theory Probab. Appl. 9, no. 4 (1964), 521-537. (Teor. Veroyat-
nosl. i Primenen. 9, no. 4 (1964), 577-594)
[3] MODE, c.J.: Mullitype branching processes, Elsevier, 1971.
VP. Chistyakov

475
BRANCHING PROCESS, AGE-DEPENDENT
Editorial comments. Additional references can be found
in the article Branching process.
AMS 1980 Subject Classification: 60J80, 58F14
BRANClDNG PROCFSS WITH A FINITE NUMBER
OF PARTICLE TYPES - A model of a branching pro-
cess which is a special case of a Markov process with a
countable set of states. The state of the branching vec-
tor is described by the random process
IJ-{/) = (PI (I), ... ,p,,(t)),
the k-th component, J.Lk(t), of which shows that at time
t there are J.Lk(t) particles of type Tk . The principal pro-
perty by which branching processes differ from Markov
processes is that the particles existing at the moment t 1
produce daughter particles at any subsequent moment
t 1 + t, t >0, in a mutually independent manner. The
generating functions
Fk(/, Xl> . . ,xn) =
= E(x),(t) , ... ,x:;-(t) IILk (0) = I; IL;(O)=O, i*k)
satisfy the system of equations
Fk(/+T,XI>'" ,xn) =
= Fk(/, FI(T, Xl> . ,xn), ... ,Fn(T, Xl> ,xn))
with initial conditions
Fk(O,XI>'" ,xn) = Xb k=I, ... ,n.
The equations (*) are satisfied by discrete-time and
continuous-time processes.
In the case of discrete time, the matrix of mathemati-
cal expectations
A(/) = II Aij(t) II,
aF I
Aij(t) = ax; x,~ . . ~x.~1
is the t-th power of the matrix A =A (1): A (t)=A I. If
A is indecomposable and aperiodic, it has a simple
positive eigen value A which is larger than the moduli
of the other eigen values. In this case, as t -HX;,
A;/t) = U;V}I/ +0(;\'),
where (u 1, . . . ,Un), (V 1, . . . ,Vn) are the right and left
eigen vectors of A which correspond to A. Branching
processes with an indecomposable matrix A are said to
be subcritical if A< I, supercritical if A< 1 and critical if
A= 1 and if at least one of the functions
F k (1, XJ, . . . ,xn ) is non-linear. The concept of critical-
ity is defined in a similar manner for continuous-time
processes.
The asymptotic properties of a branching process sig-
nificantly depend on its criticality. Subcritical and criti-
cal processes die out with probability one. The asymp-
totic formulas (as t~oo) for the probabilities,
Qk(t) = P{ILI(t)+ ... +ILn(tO
476
IILk (0) = 1; P.;(O)=O, i*k},
and the theorems on limit distributions of the number
of particles [2], are analogous to the respective results
for processes involving single-type particles (cf. Branch-
ing process). Asymptotic properties in the near-critical
case (t~oo, A~l) have been studied [3]. Processes
with a decomposable matrix of mathematical expecta-
tions have also been investigated [4].
References
[1] KOLMOGOROV, A.N. and DMlTRIEv, N.A.: 'Branching stochas-
tic processes', Dokl. Akad. Nauk SSSR 56, no. 1 (1947), 5-8
(in Russian).
[2] SEWASTJANOW, B.A. [B.A. SEVAST'YANOV]: Verzweigungspro-
zesse, Akad. Wissenschaft. DDR, 1974.
[3] CuISTYAKOV, V.P.: 'On transition phenomena in branching
stochastic processes with several types of particles', Theory Pro-
bab. Appl. 17, no. 4 (1972), 631-639. (Teor. Veroyatnost. i
Primenen. 17, no. 4 (1972), 669-678)
[4] OGURA, Y.: 'Asymptotic behaviour of multitype
Galton -- Watson processes', J. Math. Kyoto Univ. 15 (1975),
251-302.
v.P. Chistyakov
Editorial comments. Additional references can be found
in the article Branching process.
(*) AMS 1980 Subject Classification: 60J80, 60J85,
58F14
BRANClDNG PROCESS WITH A RANDOM
MEDIUM, branching process in a random medium - A
time-inhomogeneous branching process in which the
inhomogeneity is random. Let g= {go, gl, ... } be a sta-
tionary sequence of random variables (the value of gl is
interpreted as the state of the 'medium' at the moment
of time t), and let to each possible state of the medium
g correspond a probability distribution (Pk(g)} of the
number of descendants of a single particle:
Pk(~) ;;. 0, ~Pk(~) = I,
k=O
~(s) = ~Pk(~)sk.
k=O
In order to construct a trajectory {J.L(O), J.L(1), ... } of a
branching process in a random medium the value of
J.L(O) = m and the trajectory g of the states of the
medium are fixed, and J.L(t + I) is determined for each
t =0, 1, ... , as a sum of J.L(t) independent random
variables with distribution (Pk(gl)}' Such a complica-
tion of the branching Galton - Watson process is
natural if, for example, the branching process in a ran-
dom medium is regarded as a model of a biological
population.
The properties of branching processes in a random
medium are analogous to those of ordinary branching
processes. For instance, the generating distribution
function of J.L(t), under the condition J.L(O)= 1, has the
form

(for a branching Galton - Watson process, i.e. for
P{~t=O} = 1, the right-hand side of (*) equals the t-fold
iteration of Fo(s)). Branching processes in a random
medium may be subcritical, critical and supercritical:
the 'criticality parameter' here (see [1]) is the variable
00
p = EQ,ln~kPk(~) = EQ,lnF~(l)
k =0
(for ordinary branching processes, the 'criticality
parameter' is the mathematical expectation of the
number of 'descendants' of a single particle). If p<O,
the branching process in a random medium is said to
be subcritical and, for the random variable
q(f) = limP{M(t)=OIM(O)=I, I}
t~oo
which is the probability of extinction of the branchin~
process in a random medium for a given trajectory ~,
the relation
P{q(~=1} = 1
is valid. There is also the analogue of the limit theorem
of the subcritical Galton - Watson brancIEng process:
For almost all realizations of the sequence ~ the limits
lim P{M(t)=klM(O)=I, M(tO, I} = pZ(I)
1->00
exist and satisfy
00 _
~p~(~) = 1.
k=1
If p = 0, the branching process in a random medium is
said to be critical, and
P{ q(~)= 1} = 1
and, for almost all realizations of ~,
lim P{M(t)=k I M(O)= I, M(tO, ~} = O.
I~oo
If p>O, the branching process in a random medium is
said to be supercritical; in such a case
P{ q(~)< 1} = I
and, if certain additional conditions are met, there
exists for almost all ~ a non-negative random variable
w,
W = lim M(t) EW = 1.
hoo F~ (1) ... F~_l (I)'
References
[1] A THREVA, K.B. and NEV, P.: Branching processes, Springer,
1972.
A.M. Zubkov
Editorial comments. Additional references can be found
in the article Branching process.
AMS 1980 Subject Classification: 60J80
BRANClDNG PROCESS WITH DIFFUSION - A
model of a branching process in which the reproducing
particles diffuse in some domain G. Let the domain G
be r-dimensional, with an absorbing boundary 8G, and
let the particles in the domain itself execute mutally
BRANCHING PROCESS WITH IMMIGRATION
independent Brownian motions. Each particle in G is
independently converted, within a time At, into n parti-
cles with a probability of PnAt + 0 (At), n:#= 1, At ~O. Let
the daughter particles begin their independent evolution
from the point of their genesis. Let
n -:::0
be the generating function of {Pn}, let PI = - ~n7'=IPn'
and let /LxAA) be the number of particles in a set A C G
at the moment t if there initially was one particle at the
point x EO G. The generating functional
H(t;x, s( = Eexp [[lnS(Y)/Lx.t(q)]
satisfies the quasi-linear parabolic equation
;=1
aax,I[ +
2
f(H) = aH
at
with the initial condition
H(O, x, sO) = sex)
and the boundary condition
H(t, x, sO) Ix~ac = O.
Denote by 0<'\1 <'\2 ~'\3 ~ . " the eigen values, and
let 1>1 (x 0 be the eigen function of the problem
~
r f i2 +A</> _
- 0,
_
!f>(x) IHac - 0,
;=1 ax;
corresponding to ,\ I. As t ~ 00 the asymptotic relation
E/Lx,(G) ~ K e(a-AIlI</>I(X)
holds. For this reason the problem is said to be subcriti-
cal if a <'\1, critical if a ='\1 and supercritical if a >'\1.
If a ~'\I' a branching process with diffusion dies out
with probability one, while if a>'\\o both the probabil-
ity of dying out and the probability of the event
f1xAG)~oo as t~oo will in general be positive.
Depending on their criticality, branching processes with
diffusion obey limit theorems analogous to those valid
for branching processes without diffusion.
References
[1] SEWASTJANOW, B.A. [B.A. SEVAST'YANOV]: Verzweigungspro.
zesse, Akad. Wissenschaft. DDR, 1974. BAS '
. . evastyanov
Editorial comments. Additional references can be found
in the article Branching process.
AMS 1980 Subject Classification: 60J60, 60J80,
58F14
BRANClDNG PROCESS WITH IMMIGRATION - A
model of a branching process (discrete-time or
continuous-time, with one or several types of particles,
etc.) in which new particles may appear not only as the
products of division of the original particles, but also as
a result of immigration from some 'external source'.
477
BRANCHING PROCESS WITH IMMIGRATION

For instance, let

Xt,i,Yt , 1=0,1, ... ; i=I,2, ... ,
lim
t~oo
p{ In(l In +p.(/,.;:
1
~x
} = x, 0,.;:,.;:1
"",x~.

be independent random variables with generating func-

References
tions 00 [l] ZUBKOV, A.M.: 'Life-like periods of a branching process with
F(s) = ~P{Xt,i=k}Sk, immigration', Theory Probab. Appl. 17, no. 1 (1972), 174-183.
k=Q (Teor. Veroyatnost. i Primenen. 17, no. 1 (1972), 179-188)
00 [2] PAKES, A.G.: 'Further results on the critical Galton-Watson
G(s) = ~P{Yt=k}Sk, process with immigration', J. Austral. Math. Soc. 13, no. 3
k=O
(1972),277-290.
respectively; the branching Galton-Watson process [3] FOSTER, J.H. and WlLUAMSON, J.A.: 'Limit theorems for the
Galton-Watson process with time-dependent immigration', Z.
with immigration may be defined by the relations
Wahrsch. Verw. Geb. 20 (1971),227-235.
J4(0)=0, where J4(t) is the number of particles and [4] SENETA, E.: 'An explicit limit theorem for the critical
Galton - Watson process with immigration', J. Roy. Statist.
p.(1 + 1) = Xt, 1 + ... + xt,fi(t) + Yt , 1 =0, 1, ....
Soc. Ser. B 32, no. 1 (1970), 149-152.
Here, the variable Xt, i is interpreted as the number of [5] SENETA, E.: 'On the supercritical Galton - Watson process with
daughter particles of the i-th particle of the t-th genera- immigration', Math. Biosci. 7 (1970), 9-14.
[6] FOSTER, J.H.: 'A limit theorem for a branching process with
tion, while the variable Y t is interpreted as the number state-dependent immigration', Ann. of Math. Statist. 42, no. 5
of the particles which have immigrated into the t-th (1971), 1773-1776.
A.M. Zubkov
generation. The generating functions
Editorial comments. Additional references may be found
Ht(s) = E{s/l!t) I p.(0)=0} in the article Branching process.
are given by the recurrence relations AMS 1980 Subject Classification: 60J80, 60J85,
Ho(s) = 1, Ht+l(S) = G(s)Ht(F(s. 60GXX
The Markov chain J4(t) corresponding to the
BRANCIUNG PROCESSES, REGULARITY OF - A
Galton - Watson branching process with immigration is
property of branching processes ensuring that the
recurrent if EXt,; < 1 and Eln(l + Yt ) < 00 or EXt,; = 1
number of particles at any moment of time is finite.
and B=DXt,;>2C=2EYt , and is transient if EXt,; = 1
The problem of the regularity of a branching process is
and B<2C or EXt,;> 1. For the Markov chain J4(t) to
be ergodic, i.e. for the limits usually reduced to the problem of uniqueness of the
solution of some differential or integral equation. For
lim P{p.(/)=k}
t~oo
= Pk instance, in a continuous-time branching process the
to exist and to satisfy differential equation
00 aF(/; s) = f(F(t; s
~Pk = 1, at
k=O
with the initial condition F(O; s)=s has a unique solu-
it is necessary and sufficient [3] that tion if and only if, for any >0, the integral
1
j I-G(s) ds < 00.
j
1
dx
o F(s)-s
l-J(x)
This condition is met, in particular, if EXt,; < 1 and
is divergent. In the branching Bellman - Harris process
Eln(l + Y t )< 00. If EXt,; = 1, B>O, C< 00, then [4]
the generating function F(t; s) of the number of parti-
lim p{1:l!!.!.l,,;;;x}
t~oo Bt
= I
r(2CB -\)
jXy2CB-I-le-Y dy, x;;;;'O.
0
cles is the solution of the non-linear integral equation

If A = EXt,; > 1 and Eln( 1+ Yt )< 00, then there exists F(t;s) = jh(F(t-u;sdG(u)+s(I-G(I, (*)
[5] a sequence of numbers Ct~O, Ct / Ct + I ~A, such that
where G(t) is the distribution function of the lifetimes
P {lim clp.(t) exists and is POSitive} = 1. of particles and h (t) is the generating function of the
I~CC
number of daughter particles ('direct descendants') of a
single particle. If, for given to, CI, C2>0 and an integer

In branching processes with immigration in which the
immigration takes place at p.(t) = only, i.e. n ~ 1, the inequalities
p.(t + 1) = X I. l + ... + XI./l!I) +80 ,/l!t) YI, t =0,1, ... , clt n ,,;;; G(t) ~ C21"
where 0ij is the Kronecker symbol, the following rela- are valid for all O,,;;;t,,;;;t o, the solution of equation (*) is
tion is valid if EXt,i=l, l<EX?;<oo and O<EYt<oo: unique if and only if the equation

478

!.!/!..
n
= h(CP)-1
dt
with initial conditions
cf>(0) = I, cp(r)(o) = 0, r = I, ... , n -I,
has a unique solution
o ,;;; cf>(t) ,,;;;;; 1.
For a branching process described by equation (*) to
be regular, it is necessary and sufficient for the integral
f'o xl-l/n(l-h(l-x))l/n
dx
to diverge for any >0.
References
[1) SEWASTJANOW, B.A. [B.A. SEVAST'YANOV): Verzweigungspro-
zesse, Mad. Wissenschaft. DDR, 1974. ,
B.A. Sevast yanov
Editorial comments. Additional references can be found
in the article Branching process.
AMS 1980 Subject Classification: 60J60, 60GXX
BRANDT SEMI-GROUP - A semi-group with a zero
element in which to each non-zero element a there
correspond uniquely determined elements e, f, a' ES
such that ea = af = a and a' a =f, and for any two non-
zero idempotents gl. g2 ES one has gl Sg2=1=O. The ele-
ments e and f in the definition will in fact be idempo-
tent and fa'=a'e=a' and aa'=e. Moreover, in a
Brandt semi-group each one of the conditions
ac = bc=l=0, ca = cb=l=O, implies a = b, while the condi-
tions ab=l=O and bc=l=O imply abc =1=0.
The partial groupoid obtained by excluding the zero
element from a Brandt semi-group is said to be a
Brandt groupoid. The concept of a Brandt semi-group
was introduced by H. Brandt in [I], and the concept of
a Brandt groupoid was introduced in the same publica-
tion by implication. The concept of a Brandt groupoid
is an abstraction of the system of normal ideals of
semi-simple linear algebras with respect to the so-called
proper multiplication (d. [2], [3], Chapt. 6). The signi-
ficance of a Brandt semi-group in the theory of semi-
groups lies in the fact that Brandt semi-groups are just
completely O-simple inversion semi-groups (d. Inversion
semi-group; Completely-simple semi-group). A semi-
group will be a Brandt semi-group if and only if it is
isomorphic to a Rees semi-group of matrix type with a
unit sandwich matrix over a group with an added zero
element.
References
[1) BRANDT, H.: 'Ueber eine Verallgemeinerung des Gruppen-
begriffes', Math. Ann. 96 (1927), 360-366.
[2) DEURING. M.: Algebren, Springer, 1935.
(3) JACOBSON, N.: The theory oj rings, Amer. Math. Soc .. 1943.
(4) SUSHKEVICH, A.K.: The theory oJ generalized groups, Kharkov-
Kiev, 1937 (in Russian).
BRAUER GROUP
(5) CLIFFORD. A.H. and PRESTON. G.B.: The algehraic theory oj
semigroups, 1-2, Amer. Math. Soc., 1961-1967.
L.N. Shevrin
AMS 1980 Subject Classification: 20MXX
BRAUER GROUP of a field k - The group of classes
of finite-dimensional central simple algebras (cf. Cen-
tral simple algebra) over k with respect to the
equivalence defined as follows. Two central simple k-
algebras A and B of finite rank are equivalent if there
exist positive integers m and n such that the tensor pro-
ducts A kMm(k) and BkMn(k) are isomorphic k-
algebras (here Mr(k) is the algebra of square matrices
of order rover k). The tensor multiplication of alge-
bras induces an Abelian group structure on the set of
equivalence classes of finite-dimensional central simple
algebras. This group is also known as the Brauer
group of k and is denoted by Br k. The zero element of
this group is the class of full matrix algebras. while the
element inverse to the class of an algebra A is the class
of its opposite algebra. Each non-zero class contains.
up to isomorphism, exactly one division algebra over k
(i.e. a skew-field over k).
Brauer groups were defined and studied in several
publications by R. Brauer, E. Noether, A. Albert, H.
Hasse and others. starting in the nineteen-twenties (see.
for example. [6]). The most complete results, including
the computation of the Brauer group, were obtained for
numher fields in connection with the development of
class-field theory. The general form of the reciprocity
law is formulated in terms of Brauer groups.
The Brauer group is zero for any separably-closed
field and any finite field. For the field of real numbers
the Brauer group is a cyclic group of order two and its
non-zero element is the class of the quaternion algebra.
If k is the field of p-adic numbers or any locally com-
pact field that is complete with respect to a discrete
valuation, then its Brauer group is isomorphic to Q / Z,
where Q is the additive group of rational numbers and
Z is the additive group of integers. This fact is of
importance in local class-field theory.
Let k be an algebraic number field of finite degree or
a field of algebraic functions in one variable with a fin-
ite field of constants. Then there exists an exact
sequence of groups
IOV 2:
0 ..... Brk ..... :LBrk, ..... Q/Z ..... 0,
where v runs through all possible norms of the field k,
kv are the respective completions of k, and the
homomorphism inv is induced by the natural imbed-
dings k--'>kv' The image of an element from Brk in
Br k v is called a local invariant, the homomorphism ~
is the summation of local invariants. This fact is esta-
479
BRAUER GROUP
blished in global class-field theory.
If k is a field of algebraic functions in one variable
over an algebraically closed field of constants, then its
Brauer group is zero (Tsen's theorem). The case of an
arbitrary field of constants is treated in [4] and in [7].
The Brauer group depends functorially on k, i.e. if K
is an extension of the field k, a homomorphism
Br k ~ Br K is defined. Its kernel, denoted by Br(K / k),
consists of classes of algebras splitting over K.
The construction of cross products with the aid of
factor systems [5] results in a cohomological interpreta-
tion of Brauer groups. For any normal extension K / k
there exists an isomorphism
Br(K / k) ~ H 2 (K, K'),
where H 2(K, K') is the second Galois cohomology
group with coefficients in the multiplicative group K'
of K. Moreover, the group Br k is isomorphic to
H2(f, t), where k is the separable closure of k. A
central simple algebra is assigned its class in the Brauer
group by the coboundary operator
0: HI(K, PGL(n, K ~ H 2 (K, K')
in the cohomology sequence corresponding to the exact
group sequence
1 ~ K' ~ GL(n, K) ~ PGL(n, K) ~ 1,
where GL(n, K) and PGL(n, K) are the linear and the
projective matrix groups of order n X n. Here the set
H1(K, PGL(n, K)) is interpreted as the set of k-
isomorphism classes of central simple algebras of rank
n 2 over the field k which split over k, or as the set of
classes of k-isomorphic Brauer - Severi varieties of
dimension n -1, possessing a point that is rational over
K (d. Brauer-Severi variety).
All Brauer groups are periodic groups. The order of
any of its elements is a divisor of the number n, where
n 2 is the rank of the skew-field representing this ele-
ment.
The cohomological interpretation of the Brauer
group makes it possible to consider it as the group of
classes of ex~ensions of the Ga~!s group of the separ-
able closure k / k by the group k .
A generalization of the concept of a Brauer group is
the Brauer - Grothendieck group, whose definition is
analogous to that of the Brauer group, except that the
central simple algebras are replaced by Azumaya alge-
bras [7].
References
[1] CASSELS. J.W.S. and FROLICH, A. (EDS.): Algebraic number
theon', Acad. Press, 1967.
[2] BOURBAKI. N.: Elements of mathematics. AIKebra: AlgebraiC
structures. Linear algebra. 1, Addison-Wesley. 1974,
Chapt.l;2 (translated from the French).
[3] SERRE. J.-P.: Cohomologie Galoisienne, Springer, 1964.
[4] FADDEEV, O.K.: 'On the theory of algehras over fields of alge-
hraic functions in one variahle', Vestnik LeninKrad. Univ .. no. 7
480
(1957),45-51 (in Russian), English summary.
[5] CHEBOTAREV, N.G.: Introduction to the theory of algebra... ,
Moscow-Leningrad, 1949 (in Russian).
[6] OEURING, M.: Algebren, Springer, 1935.
[7] GROTHENDIECK, A.: 'Le groupe de Brauer I, II, III', in Dix
exposes sur la cohomologie des schemas, North-Holland, 1968,
pp.46-188.
[8] MILNE, l.S.: Etale cohomology, Princeton Univ. Press, 1980.
V.A. Iskovskikh
Editorial comments, For the construction of cross pro-
ducts with the aid of factor systems see also Cross pro-
duct; Extension of a group. The latter article contains the
notion of a system of factors.
A recent result in the theory of Brauer groups is a
theorem of Merkuryev and Suslin [A1], which in its simplest
form asserts that Br k is generated by the classes of alge-
bras that split over a cyclic extension of k, provided that k is
a field of characteristic zero containing all roots of unity.
The proof is based on the close relationship between the
theory of Brauer groups and algebraic K-theory. An alge-
bra A over a commutative ring is an Azumaya algebra if it is
finitely generated and central over k and separable.
References
[A1] SUSLIN, A.: 'Plenary adress', in Proc. Internat. Congress of
Mathematicians Berkely, 1986, Forthcoming.
AMS 1980 Subject Classification: 13A20, 16A 16
BRAUER - SEVERI VARIETY An algebraic
variety over a field k that, if considered over the alge-
braic closure k of k, becomes isomorphic to a projective
space,
The arithmetic properties of such varieties were stu-
died in 1932 by F. Severi; F. Chiitelet subsequently
discovered a connection between Brauer - Severi
varieties and central simple algebras (d, Central simple
algebra) over k and the Brauer group.
The simplest non-trivial example of a one-
dimensional Brauer - Severi variety is the projective
conic section Q:
X5 +XT +x~ = 0
on the real projective plane p~, Over the field of com-
plex numbers C this variety is isomorphic to the projec-
tive line Pt. The set of all one-dimensional
Brauer - Severi varieties, considered up to isomorphism,
is in a one-to-one correspondence with the set of pro-
jective non-degenerate conic sections (considered up to
projective equivalence over k). which is in turn in a
one-to-one correspondence with the set of non-
isomorphic generalized quaternion algebras over k. In
the above example the conical section Q corresponds to
the algebra of ordinary quaternions.
In the more-dimensional case, the set of classes of n-
dimensional Brauer - Severi varieties (i.e.
Brauer- Severi vanehes distinguished up to k-
isomorphism) may be identified with the Galois coho-

mology group HI(k, PGL(n + 1, k where
PGL(n + 1, k)
is the projective group of automorphisms
of the projective space PI: [3), [4). This cohomology
group describes the classes of k-isomorphic central sim-
ple k-algebras of rank (n + 1)2 (i.e. forms of the matrix
algebra Mn + 1(k)). The connection between
Brauer- Severi varieties and central simple algebras is
more explicitly described as follows. To a k-algebra A
of rank r2 one associates the variety X of its left ideals
of rank r, which is defined as a closed subvariety of the
Grassmann manifold of all k-linear subspaces of dimen-
sion r in A. In certain cases the variety X may be
defined by norm equations - e.g. in the case of
quaternion algebras. The connection between
Brauer - Severi varieties and algebras is taken advan-
tage of in the study of both [1), [4).
The most significant properties of Brauer- Severi
varieties are the following. A Brauer - Severi variety is
k-isomorphic to a projective space PI: if and only if it
has a point in the field k. All Brauer - Severi varieties
have a point in some finite separable extension K of k
[ 1).
The Hasse principle applies to Brauer - Severi
varieties defined over an algebraic number field.
The field of rational functions k (X) on a
Brauer - Severi variety X is the splitting field of the
corresponding algebra A; moreover, an arbitrary exten-
sion K of k is the splitting field for A if and only if X
has a K-point [4).
In the context of the generalization of the concepts
of a central simple algebra and the Brauer group to
include schemes, the Brauer- Severi varieties were gen-
eralized to the concept of Brauer - Severi schemes [2).
Let 1: p--,>x be a morphism of schemes. A scheme P is
called a Brauer-Severi scheme if it is locally iso-
morphic to a projective space P1 over X in the etale
topology of X. A scheme P over a scheme X is a
Brauer-Severi scheme if and only if 1: P--,>X is a
finitely-presented proper flat morphism and if all of its
geometrical fibres are isomorphic to projective spaces
[2).
References
[I) CHA":ELET, F.: 'Variations sur un theme de H. Poincare', Ann.
Sci. Ecole Norm. Sup. (3) 61 (1944), 249-300.
[2] GROTHENDIECK, A.: 'Le groupe de Brauer', in Dix exposes sur
la cohomologie des schemas, North-Holland, 1968, pp. 1-21.
[3] SERRE, J.-P.: Cohomologie Galoisienne, Springer. 1964.
[4] ROGUEITE, P.: 'On the Galois cohomology of the projective
linear group and its applications to the construction of generic
splitting fields of algebras', Math. Ann. 150 (1963), 411-439.
VA. Iskovskikh
Editorial comments. Thus a Brauer-Severi variety of
dimension n is a k / k-form of P~.
AMS 1980 Subject Classification: 14M99, 14GXX
BRIOT - BOUQUET EQUATION
BREAKING POINT, angle point - A singular point of
a plane curve with the property that two branches of
the curve end in it in such a way that each has a (one-
sided) tangent at the point different from the other. For
example, the origin is a breaking point of the curve
y=x /(l+e 1 / x ) (see Fig.). The left and right deriva-
tives are different at a breaking point.
y
x
A.B. Ivanov
AMS 1980 Subject Classification: 53A04
BRIANCHON THEOREM - In any hexagon (Fig.)
circumscribed around a curve of the second order (a
Brianchon hexagon) the straight lines connecting the
opposite corners of the hexagon intersect at one point
(Brian chon 's point). This theorem is dual to the Pascal
theorem, and was demonstrated by Ch.l. Brianchon in
1806.
A.B. Ivanov
AMS 1980 Subject Classification: 51 N1 0
BRIOT- BOUQUET EQUATION - An ordinary dif-
ferential equation
(I)
where m is a positive integer and the function 1 is ana-
lytic at x=y=O, h(O,O)=f:O, 1(0,0)=0. It was shown
by C. Briot and T. Bouquet [1) that any equation of the
type
a(z, w)w' = (3(z, w),
where a(O, 0) = /3(0, 0) = and a and /3 are analytic at
the origin, can be reduced, by means of a special local
changes of the variables, to a finite number of equa-
tions of type (1). Equation (1) always (except for the
case where m = 1 and h(O, 0) is a natural number) has
a unique solution in the form of a formal power series:
=
Y = ~(x) ~IX+~2X2+ " ' , (2)
which converges for sufficiently small I x I if m = 1,
and can diverge for all x=f:0 if m> 1. In (1), let
f fo(x)+ fl(x)y,
then, for the series (2) to converge, it is necessary and
481
BRIOT- BOUQUET EQUATION
sufficient to meet m -I conditions concerning the coef-
ficients of the Taylor series of /0 and /1; all the coeffi-
cients are included in these conditions, so that the
existence or non-existence of an analytic solution
y = ~(x) of equation (1) cannot be proved by any par-
tial sum of the Taylor series of / (cf. [2], [3]). For the
case of a general function / there are
(m -I)+(m -I) X 00 such conditions, [4]. Accordingly,
the Briot - Bouquet equation is sometimes referred to
as equation (1) with m>l.
References
[I] BRIOT, C. and BOUQUET, T.: 'Recherches sur les proprietes des
equations differentielles', J. Eco/e Poly technique 21, no. 36
(1856), \33-198.
[2] BIEBERBACH, L.: Theorie der gewohnlichen Differentia/-
g/eichungen auffunktionentheoretischer Grund/age dargestellt,
Springer, 1965.
[3] BRYUNO, A.D.: 'Analytical form of differential equations.
Introduction', Trans. Moscow Math. Soc. 25 (1971), 134-15J.
(Trudy Moskov. Mat. Obshch. 25 (1971), 120-138)
[4] MARTINET, J. and RAM IS, J.P.: 'Problemes de modules pour
des equations differentielles du premier ordre', Pub/. Math.
IHES 55 (1982), 63-164.
A.D. Bryuno
AMS 1980 Subject Classification: 34A30
BROUWER LATTICE, Brouwer structure, Brouwer
algebra - A distributive lattice in which for each pair of
elements a, b there exists an element, called the
pseudo-difference (frequently denoted by a *- b), which is
the smallest element c possessing the property b +c;;;;'a.
An equivalent description of a Brouwer lattice is as a
variety of universal algebras (cf. Universal algebra) with
three binary operations u,
nand *-, which satisfies
certain axioms. The term 'Brouwer algebra' was intro-
duced in recognition of the connection between
Brouwer lattices and Brouwer's intuitionistic logic.
Instead of Brouwer lattices the so-called pseudo-Boolean
algebras are often employed, the theory of which is
dual to that of Brouwer lattices. Any Brollwer lattice
can be converted to a pseudo-Boolean alge) fa by the
introduction of a new order (a~'b)=>(b~a), and of
new unions and intersections according to the formulas
(a U .h) = (a n h), (a n' h) = (a U h)
and the operation of relative pseudo-complementation
a=>b which corresponds to the pseudo-difference a*- h.
Conversely, any pseudo-Boolean algebra can be
regarded as a Brouwer lattice. The term 'Brouwer lat-
tice' is sometimes used to denote a pseudo-Boolean
algebra (see, for instance, [2]).
References
[IJ McKINSEY. 1.C.C. and TARSKI. A.: 'The algehra of topology'.
Ann. oj Mmh. (2) 45. Ill). 1(1944). 141191.
[21 BIRKIIOFI'. G.: LallieI' lhl'orr. Amer. Math. Soc:. 1967.
V.A. Yankov
Editorial comments. In Western literature pseudo-
4X2
Boolean algebras are more commonly called Heyting alge-
bras. Complete Heyting algebras (often called frames or
locales) have been extensively studied on account of their
connections with topology: the lattice of open sets of any
topological space is a locale, and locales can in some
respects be considered as generalized topological spaces.
See [A1], [A2], [A3].
References
[A1] FOURMAN, M.P. and SCOTT, D.S.: 'Sheaves and logic', in
Applications of sheaves, Lecture notes in Math., Vol. 753,
Springer, 1979, pp. 302-401.
[A2] JOHNSTONE, P.T.: Stone spaces, Cambridge Univ. Press,
1983.
[A3] SIMMONS, H.: 'A framework for topology', in Logic colloquium
'77, Studies in logic and foundations of math., Vol. 96,
North-Holland, 1978, pp. 239-251.
AMS 1980 Subject Classification: 06D20, 03G10
BROUWER THEOREM - 1) Brouwer's fixed-point
theorem: Under a continuous mapping j: S ~S of an
n-dimensional simplex into itself there exists at least
one point XES such that j(x)=x; this theorem was
proved by L.E.1. Brouwer [I]. An equivalent theorem
had been proved by P.G. Bohl [2] at a somewhat earlier
date. Brouwer's theorem can be extended to continuous
mappings of closed convex bodies in an n-dimensional
topological vector space and is extensively employed in
proofs of theorems on the existence of solutions of vari-
ous equations. Brouwer's theorem can be generalized to
infinite-dimensional topological vector spaces.
References
[I] BROUWER, L.E.].: 'Ueber eineindeutige. stetige Transforma-
tionen von FUichen in sich'. Math. Ann. 69 (1910), 176-180.
[21 BOHL, P.: 'Ueber die Beweging eines mechanischen Systems in
der Nahe einer Gleichgewicht>lage'. J. Reine Angew. Math. 127
(1904), 179-276.
VI. Sobolev
Editorial comments. There are many different proofs of
the Brouwer fixed-point theorem. The shortest and concep-
tually easiest, however, use algebraic topology.
Completely-elementary proofs also exist. Cf. e.g. [A 1],
Chapt. 4. In 1886, H. Poincare proved a fixed-point result
on continuous mappings f: P--,>P which is now known to
be equivalent to the Brouwer fixed-point theorem, [A2].
There are effective ways to calculate (approximate) Brouwer
fixed pOints and these techniques are important in a multi-
tude of applications including the calculation of econom'c
equilibria, [A1]. The first such algorithm was proposed by H.
Scarf, [A3]. Such algorithms later developed in the so-called
homotopy or continuation methods for calculating zeros of
functions, cf. e.g. [A4], [A5].
References
[A1] ISrRATI'sCC. V.1. Fixedpointtheory. Reidel. 1981.
[A2] POINCAR(. H.: 'Sur les courbes definies par les equations
differentielles', J de Math. 2 (1886).
[A3J S( ARI. II .. 'The approximation of fixed points of continuous
mappings', SIAM J Apo/. Math. 15 (1967),1328-1343.
[A4] KARAMADlAN, S. (ED.): Fixed pOints. Algorithms and applica-
tions, Acad. Press, 1977.

[A5] ALLGOWER, E. and GEORG, K.: 'Simplicial and continuation
methods for approximating fixed pOints and solutions to sys-
tems of equations', SIAM Rev. 22 (1980), 28-85.
2) Brouwer's theorem on the in variance of domain:
Under any homeomorphic mapping of a subset A of a
Euclidean space En into a subset B of that space any
interior point of A (with respect to En) is mapped to
an interior point of B (with respect to En), and any
non-interior point is mapped to a non-interior point. It
was proved by L.E.l. Brouwer [1].
References
[I] BROUWER, L.E.1.: 'Ueber Abbildungen von Mannigfaltig-
keiten', Math. Ann. 71 (1912), 97-115.
M.J. VoTtsekhovskiT
Editorial comments. For a modern account of the
Brouwer invariance-of-domain theorem cf. [A 1]. Chapt. 7,
Sect. 3. The result is Important for the idea of the
topological dimension (dim e = n).
References
[A 1] DUGUNDJI. J.: Topology, Allyn and Bacon, 1966.
AMS 1980 Subject Classification: 54H25, 55M20,
57A15
BROWNIAN MOTION - The process of chaotic dis-
placements of small particles suspended in a liquid or
in a gas which is the result of collisions with the
molecules of the medium. There exist several
mathematical models of this motion [1]. The model of
Brownian motion which is the most important one in
the theory of random processes is the so-called Wiener
process, and the concept of Brownian motion is in fact
often identified with this model.
References
[1] PAVLOV, V.P.: 'Brownian motion', in Large Soviet Encyclo-
pedia, Vol. 4 (in Russian).
Editorial comments. See also Wiener measure.
References
[A 1] ITt>, K. and McKEA"', H.P .. JR.: Diffusion processes and their
sample paths. Springer. 1974, Chap!. 1; 2.
AMS 1980 Subject Classification: 60J65, 58G32
BRUHAT DECOMPOSITION - A representation of a
connected algebraic reductive group G as the union of
double cosets of a Borel subgroup, parametrized by the
Weyl group of G. More exactly, let B, B - be opposite
Borel subgroups of a reductive group G; let V, V - be
the respective unipotent parts of B, B - (d. Linear
algebraic group) and let W be the Weyl group of G. In
what follows w denotes both an element of Wand its
representative in the normalizer of the torus B B- , n sieve to prove that all large even numbers N can be
since the construction presented below is independent
of the representative chosen. The group
V;' = V n w V - w - I will then be considered for each
WE' W. The group G is then representable as the union
of the non-intersecting double cosets BwB (w E' W), and
BRUN SIEVE
the morphism V;' XB~BwB x,Y)l->xwy) is an iso-
morphism of algebraic varieties. An even more precise
formulation of the Bruhat decomposition will yield a
cellular decomposition of the projective variety G / B.
Namely, if Xo is a fixed (with respect to the left shifts
by elements from B) point of G / B (such a point
always exists, cf. Borel fixed-point theorem), G / B will
be the union of non-intersecting V-orbits of the type
V(w(xo, WE' W (d. Algebraic group of transforma-
tions), and the morphism V;' ~Vw(xo) (UI->u(w(xo))) is
an isomorphism of algebraic varieties. All groups V;"
being varieties, are isomorphic to an affine space; if the
ground field is the field of complex numbers, then each
of the above V-orbits is a cell in the sense of algebraic
topology so that the homology of G / B can be calcu-
lated. The existence of a Bruhat decomposition for a
number of classical groups was established in 1956 by
F. Bruhat, and was proved in the general case by C.
Chevalley [3]. A. Borel and 1. Tits generalized the con-
struction of Bruhat decompositions to the groups Gk of
k-points of a k-defined algebraic group [2], the role of
Borel subgroups being played by minimal parabolic k-
subgroups, the role of the groups V by their unipotent
radicals; the Weyl k-group W k or the relative Weyl
group was considered instead of W.
References
[I] BOREL, A.: Linear algebraic groups, Benjamin, 1969.
[2] BOREL, A. and TITS. J.: 'Groupes reductifs'. Publ. Math. IHES
27 (1965), 55-150.
[3] CHEVALLEY, c.: Classification des groupes de Lie algebriques, 2,
Paris. 1958.
VP. Platonov
AMS 1980 Subject Classification: 20G05
BRUN SIEVE - A sieve method in elementary
number theory, proposed by V. Brun [1]; it is an exten-
sion of the sieve of Eratosthenes (d. Eratosthenes,
sieve ot). The method of Brun's sieve may be described
as follows. Out of a sequence of natural numbers an:%;x
the numbers with small prime divisors are eliminated
('sieved out') leaving behind prime and almost-prime
numbers (cf. Almost-prime number) with only large
prime divisors. Let P(x) be the amount of these
numbers. It can be shown that P (x) is included
between two sums with a relatively-small number of
summands, which may be estimated from above and
from below. Thus, it is possible to evaluate from above
the number of twins in a given interval. Brun's sieve is
employed in additive number theory. Brun used his
represented in the form N = P I + P 2, where P I and P 2
contain at most 9 prime factors.
References
[1] BRUN, V.: 'Le crible d'Eratosthene et Ie theoreme de Gold-
bach', CR. Acad. Sci. Paris Ser. I Math. 168. no. 11 (1919).
544-546.
483
BRUN SIEVE
[21 GEL'FOND, A.D. and LINNIK, Yu.V.: Elementary methods in the
analytic theory of numbers, M.LT., 1966 (translated from the
Russian).
[31 TROST, E.: Primzahlen, Birkhauser, 1953.
N.!. Klimov
Editorial comments.
References
[A11 HALBERSTAM. H. and RICHERT. H.E.: Sieve methods. Acad.
Press. 1974.
AMS 1980 Subject Classification: 10H30
BRUN lHEOREM on prime twins - The series
~ 1/ p is convergent if p runs through all (the first
members of all) prime twins. This means that even if
the number of prime twins is infinitely large, they are
still located in the natural sequence rather sparsely.
This theorem was demonstrated by V. Brun [1]. The
convergence of a similar series for generalized twins
was proved at a later date.
References
[II BRUN. v.: 'La serie 1/5 + I /7+ ... ou les denominateurs
sont 'nombres premiers jumeaux' et convergente ou finie'. Bull.
Sci. Math. (2) 43 (1919).100-104; 124-128.
[21 TROST. E.: Primzahlen, Birkhauser. 1953.
N.l. Klimov
Editorial comments.
References
[A1] HALBERSTAM. H. and RICHERT. H.E.: Sieve methods, Acad.
Press. 1974.
AMS 1980 Subject Classification: 10J15
BRUNN - MINKOWSKI lHEOREM - Let K 0 and
K I be convex sets in an n-dimensional Euclidean
space; let K", AE[O, 1], be the set of points which
divide segments with end points at any points of the
sets Ko and KI in the ratio A/(l-A) (a linear combi-
nation of Ko and K I ); and let V(A) be the n-th power
root of the volume of the set K". Then V(A) is a con-
cave function of A, i.e. the inequality
V(A[(l-p)+A2P) ;;;. (l-p)V(Ad+pV():)
is valid for all AI, A2, pE[O, 1]. The function V(A) is
linear (and the inequality then becomes an equality) if
and only if Ko and Klare homothetic. The
Brunn - Minkowski theorem can be generalized to
linear combinations of several convex sets. I t is used to
solve extremal and uniqueness problems. It was
discovered by H. Brunn in 1887, and completed and
rendered more precise in 1897 by H. Minkowski.
References
[II BUSEMAN. H.: Convex surfaces. Interscience. 1958.
[2] HADWIGER. H.: Vorlesungen fiber Inhalt, Oberjlache und Isoper
imetrie, Springer. 1957.
M.I. Vottsekhovskit
Editorial comments.
References
[A 1] LEICHTWEISS, K.: Konvexe Mengen, Springer, 1979.
AMS 1980 Subject Classification: 52A40, 52A20
484
BUBNOV - GALERKIN ME1HOD - See Galerkin
method.
AMS 1980 Subject Classification: 49015, 65N30
BUDAN - FOURIER lHEOREM - The number of
roots of an algebraic equation
!(x) =0
comprised in an interval (a, h), a <h, is equal to or is
smaller, by an even number, than 7'=tl -t2, where tl
is the number of changes in sign in the series of deriva-
tives of the polynomial f (x) at the point a, i.e. in the
series
!(a), I (a), ... ,/")(a),
while t 2 is the number of changes in sign in this series
at the point h. Each multiple root is counted according
to its multiplicity. Established by F. Budan (1822) and
1. Fourier (1820).
References
[I) Encyclopaedia of elementary mathematics, Vol. 2 Algebra.
Moscow-Leningrad. 1951, p. 331 (in Russian).
o.A. lvanova
Editorial comments. An application of the
Budan - Fourier theorem in numerical analysis may be
found in [A1], where it is used in the interpolation by spline
functions.
References
[A1] BOOR, C. DE and SCHOENBERG, U.: 'Cardinal interpolation
and spline functions VIII. The Budan-Fourier theorem for
splines and applications.', in K. Bohmer, G. Meinardus and
W. Schemp (eds.): Spline functions, Lec!. Notes in Math.,
Vol. 501, Springer, 1976.
[A2] HOUSEHOLDER, A.S.: 'Unique triangularization of a nonsym-
metric matrix', J. Assoc. Compo Mach. 5 (1958), 339-342.
AMS 1980 Subject Classification: 65H05
BUFFON PROBLEM of the needle- A classical prob-
lem in the theory of geometric probabilities, which is
rightly considered to be the starting point in the
development of this theory. First posed by G. Buffon
in 1733 and reproduced together with its solution in [1].
Buffon considered the following situation: A needle of
length 2r, where 2r<a, is thrown at random on a plane
ruled by parallel lines at a distance a from one another.
What is the probability that the needle thus thrown will
come to rest across one of the lines? Clearly, the loca-
tion of the needle is defined by the distance x from its
centre to the nearest straight line and by the acute
angle () between the needle and the perpendicular to
the line. The magnitude x lies between zero and a /2,
while () lies between zero and 'TT / 2. It is assumed that
the point (x, () is uniformly distributed in the appropri-
ate rectangle (this is equivalent to assuming that the
random variables x and () are independent and are uni-
formly distributed over (0, a /2) and (0. 'TT / 2.
Then
the sought probability is defined as the ratio between
the areas corresponding to the favourable and to all
 BURKILL INTEGRAL

possible results, and is References

[I) BOUNJAKOWSKY, W. [V.YA. BUNYAKOVSKJi): 'Sur quelques
inegalites concernant les integrales aux differences finis', Mem.
p = --
1
~!!...
f rcosfJdfJ = -.
w/2

0
4r
aw
(*)
22
In the past Buffon's problem served as the tool for the
experimental check on the Bernoulli theorem. In fact, if
the needle is thrown n times, and crosses one of the
lines in m cases, the frequency m / n should, for large
values of n, according to Bernoulli's theorem, be close
to the probability (*). This idea was used by many
workers for the determination of the number '1T by ran-
dom trials [1], [2]. Buffon also considered other similar
problems, in particular the problem of the needle cross-
ing the lines belonging to two mutually-perpendicular
systems which subdivide the plane into rectangles hav-
ing sides of lengths a and b respectively. The answer
given by Buffon to this problem is incorrect. The
correct solution
4r(a+b)-4r 2
wab
was found by P. Laplace in 1812.
References
[I) BUFFON, G.: Essai d'arithmetique morale. Supplementa
'/'Histoire Morale', 4, 1777.
[2) USPENSKY, J.V.: Introduction to mathematical probability,
McGraw-Hili, 1937.
[3) KENDALL, M.G. and MORAN, PAP.: Geometric probability,
Griffin, 1963.
A. V. Prokhorov
AMS 1980 Subject Classification: 60005
BUNDLE
Editorial comments. In the sense of a (two-dimensional)
family of objects the usual name is net. (See also Web of
spheres.)
See also Vector bundle; Fibre space; Fibration.
AMS 1980 Subject Classification: 55RXX 14F05,
32LXX, 57Q50, 57R22, 53AXX, 14C21, 14C20,
51N20
BUNYAKOVSKII INEQUALITY - An equality in
mathematical analysis, established by V.Ya.
BunyakovskiI [1] for square-integrable functions f and
g:
This inequality is analogous to Cauchy's algebraic ine-
quality
(a,b, + ... +anbnf .;;;
.;;; (at + ... +a~)(bt + ... +b~).
The BunyakovskiI inequality is also known as the
Schwarz inequality; however, BunyakovskiI published
his study as early as 1859, whereas in H.A. Schwarz'
work this inequality appeared as late as 1884 (without
any reference to the work of Bunyakovskll,,).
A cad. Sci. St. Petersbourg (7) 1 (1859), 9.
v.I. Bityutskov
Editorial comments. In Western literature this inequality
is often called the Cauchy inequality, or the
Cauchy- Schwarz inequality. Its generalization to a function
fin Lp and a function gin Lq,.1 / p+ 1 / q= 1, is called the
Holder inequality.
Cauchy's algebraic inequality stated above holds for real
numbers ai, bi , i= 1, ... ,n. For complex numbers ai, bi ,
i = 1, ... ,n, it reads
I alti; + ... +anbn 12 =
I a1 + ... +a~ I . I ~ + ... +&n I
It has a generalization analogous to the Holder inequality.
References
[A1) RUDIN, W.: Principles of mathematical analysis, McGraw-Hili,
1953.
AMS 1980 Subject Classification: 26010
BURKILL INTEGRAL - A concept which was intro-
duced by J.e. Burkill [1] to determine surface areas.
The Burkill integral is introduced in its modem form
for the integration of a non-additive function F(l) over
an n-dimensional segment (a block). Let R be a set that
can be represented as a sum (union) of a finite number
of segments (such a set is called a figure). Each
representation R = U 1 k is called a subdivision C!f the
figure R. The upper and the lower Burkill integral of the
segment function F(l) over the figure R are, respec-
tively, the upper and lower limits of the sums L.kF(ld
for all possible subdivisions as the maximum of the
diameters of the segments involved in the subdivision
tends to zero. If these integrals are equal, their common
value is the Burkill integral of F over R and is denoted
by f RF. If F is integrable over R, then F is integrable
on each figure R 1 cR. This enables one to introduce an
indefinite Burkill integral, which is an additive set func-
tion. If F is continuous, the indefinite Burkill integral is
continuous as well.
The concept of the Burkill integral can be general-
ized to include the case of a set function defined on
some class of subsets of an abstract measure space.
This class must meet a number of requirements; in par-
ticular, each set of the class must permit a subdivision
into sets also of this class that have a measure as small
as one pleases. The Burkill integral can then be defined
for any set in the class in analogy with the n-
dimensional case, the respective limits being taken as
the maximum of the measures of the constituent sets
tends to zero. The Burkill integral can be naturally gen-
eralized to set functions with values in a commutative
topological group. The Burkill integral is less general
than the subsequently introduced Kolmogorov integral,

485
BURKILL INTEGRAL
which is also known as the Burkill- Kolmogorov
integral. Any function that is Burkin-integrable is also
Kolmogorov-integrable after a suitable ordering of the
subdivisions. The converse statement is true only if cer-
tain additional conditions are satisfied. The Burkin
integral is used in constructing the Denjoy integral in
different spaces.
The name of Burkin integral is also given to a
number of generalizations of the Perron integral (AP-,
CP-, SCP-integrals), which were also introduced by
Burkill. Instead of ordinary derivatives certain general-
ized derivatives are used in the definition of these
integrals. These Burkill integrals are used in the theory
of trigonometric series.
References
[IA) BURKILL, J.C: 'Functions of intervals', Proc. London Math.
Soc. (2) 22 (1924), 275-310.
[IB) BURKILL, J.C: 'The expression of area as an integral', Proc.
London Math. Soc. (2) 22 (1924),311-336.
[2) SAKS, S.: Theory of the integral, Hafner, 1952 (translated from
the Polish).
[3) ROMANOVSKli, P.I.: 'l'Integrale de Denjoy en espaces
abstractes', Mat. Sb. 9 (51), no. I (1941),67-120.
[4) BURKILL, J.C: 'Integrals and trigonometric series', Proc. Lon-
don Math. Soc. (3) I (1951), 46-57.
AMS 1980 Subject Classification:
49F25
BVRMANN- LAGRANGE SERIES, Lagrange series
- A power series which offers a complete solution to the
problem of local inversion of holomorphic functions. In
fact, let a function w = g(z) of the complex variable z
be regular in a neighbourhood of the point z = a. and
let g(a)=b and g'(a)=f=O. Then there exists a regular
function z =h(w) in some neighbourhood of the point
w = h of the w-plane which is the inverse to g (.::) and is
such that h (b)=a. Moreover, if fez) is any regular
function in a neighbourhood of the point .:: = a, then
the composite function F(w)= /[h (11')1 can be
expanded in a Burmann - Lagranfie series ) 1 a neigh-
bourhood of the point w = h
F(w) = f\a)-l
+
,,2:x
1111
-I ,r - 1 l" (-) l z - a
d" -
ldz" 1 I - ~(z)-h
1" 11
jJ
The inverse of the function \\'=g(.::) is obtained hy set-
ting/(.::) .::.
The expansion (*) follows from Biirm{lnn's theorem
[1]: Under the assumptions made ahove on the holo-
morphic functions fi (z) and 1(.::). the latter function
may he represented in a certain domain in the .::-planc
containing a in the form
Ii z) -c Ii (/ ) -l
j "~llg(z)
Ii 1 n.
1 h]" I d"
I-dz-" -I II
1 r"
(z)
I z- (i
-~(_--)=h J II
4~6
where
R = -I-If [g(Z)-b ]m-I (t)g'(z)dt dz
m 27Ti a y g(t)-b g(t)-g(z)'
Here y is a contour in the t-plane which encloses the
r
points a and z, and is such that if is any point inside
y, then the equation g(t)=gm has no roots on y or
inside yother than the simple root t = r.
The expansion (*) for the case b =0 was obtained by
J.L. Lagrange [2].
If the derivative g'(t) has a zero of order r-1 at the
point z = a, there is the following generalization of the
Biirmann - Lagrange series for the multi-valued inverse
function [3]:
F(w)= /(a)+
+1l~1~ {::,~II [I(Z) [g~Z~~b r] t~(J(W -h)"
r
Another generalization (see, for example. [4]) refers to
functions g(z) regular in an annulus; instead of the
series (*), one obtains a series with positive and nega-
tive powers of the difference w - h.
References
VA. Skvortsov [I) BORMANN. H.: Mem. In.H. Nat. Sci, Arts. Sci. Math. Ph)'s. 2
26A39, 28A 10, (1799). 13-17.
[2A) LAGRANGE, J.L.: Mem. Acad. R. Sci. et Be/les-Iellres Berlin 24
(1770).
[2B] LAGRANGE. J.L.: 'AdJitions au memoire sur la resolution des
equations numcriques', in Oeuvres. Vol. 2. Georg Olms. 1973.
pp. 579-652.
[3] HURWITZ. A. and COCRA:-;T. R.: Vorlesungen iiher allgemeine
Funkionenenrheorie und elliptische Funktionen. 1. Springer.
1968. Chapt. 7.
[4] WHITTAKER. E.T. and WATSON. G.N,: A course oj'modern
anah sis. Cambridge U niv. Press. 1952.
[5] MARKUSHFVICll. A.1.: Theon' olIul1Clions oj' a complex variahle.
1. Chelsea. reprint. 1977 (translated from the Russian).
E. D. Solomentsev
Editorial comments. There is an exhaustive treatment of
the Lagrange - Burmann theorem and series in [A 1].
References
[A 1] HENRIe!. p, Applied and computational complex analysis. 1,
Wiley, 1974
(*) AMS 1980 Subject Classification 30810
(1\' -- h)".
BURNSIDE PROBLEM - 1) The Burnside problem
for finite groups: Do unsolvable finite groups of odd
order exist? Or, in another formulation: Ale all finite
simple non-Abelian groups of even order? This problem
is connected with the name of W. Burnside, who noted
in 1897 that all simple non-Abelian groups which were
known at that time were of even order [1]. The problem
was solved in 1962 by W. Feit and J.G. Thompson [2],
who showed that all finite groups of odd order are solv-
able.
l' 1 1
r R",.
References
[I] BUR~SIDF. W.: Theon' ojgroups oj' finite order. Camhridge
Univ. Press. 1897.

[2] FElT, W. and THOMPSON, J.G.: 'Stability of groups of odd
order', Pacific J. Math. 13 (1964), 775-1029. VD M
..
Editorial comments. In [1], Burnside also showed that
all groups of order p8 q b, where p, q are prime numbers and
a, b~O, are solvable.
2) The Burnside problem for periodic groups: This
problem was posed by W. Burnside in 1902 [1): Is a
finitely-generated group, each element of which is of
finite order, always finite (the unbounded Burnside prob-
lem)? This problem may also be formulated as fol-
lows: Are all periodic groups locally finite (cf. Locally-
finite group)? Burnside himself stressed an important
special case of this problem, when the orders of all
elements of the groups are bounded uniformly (the
bounded Burnside problem), i.e. for a certain natural
number n the identity xn = I is valid in the group. It is
the bounded variant of the Burnside problem which has
received the most attention. In other words, the object
studied was the quotient group B(d, n)=F / P of a
free group F with d~2 generators by the smallest nor-
mal subgroup P containing the n-th powers F of all
elements fEF. The following results are known:
B (d, 2) is an elementary Abelian group of order 2d;
B (d, 3) is a finite group of order 3md , where
[1), [2); B(d, 4) is a finite group (of order 212 if d=2
and of order 269 if d=3 [1), [3), [4)); B(d, 6) is a finite
group of order 2s 31 , where s = 1+(d -I)3 md , t =mr and
r=I+(d-I)2d [5), [6). The negative solution to the
bounded Burnside problem was announced in 1959 [7).
The negative solution to the unbounded Burnside prob-
lem was published in 1964 [8). Another construction of
a periodic group which is not locally finite was subse-
quently given [9). P.S. Novikov and S.1. Adyan [10)
proved in 1968 that the group B(d, n), d~2, is infinite
for all odd n~4381 (a negative solution of the bounded
Burnside problem). It was subsequently shown that the
word problem and the conjugation problem for these
values of d, n are solvable in B(d, n); B(d, n) cannot
be presented by means of a finite number of defining
relationships; all finite subgroups in B (d, n) are
Abelian, and all Abelian subgroups are cyclic; B (d, n),
d~n, do not meet the maximum or minimum condi-
tions for normal subgroups; the B(d, n), d>2, are iso-
morphically imbeddable in the group B (2, n). For a
study of the properties of the infinite groups B (d, n),
based on the advanced methods of [10), see the mono-
graph [11) in which, in particular, the bound on the
odd values of n indicated above is reduced to n ~665.
To give a final bound for the indices n to which
correspond finite or infinite groups B (d, n) is a difficult
problem. The values n = 5, 12, and n = 2m , m ~ 3, are
especially interesting in this context.
BURNSIDE PROBLEM
In accordance with the idea which has been gradu-
ally evolving ever since the middle of the nineteen-
azurov
thirties, the answer to the following question is of
importance in the theory of finite groups: Is the order
of any finite group with d generators, obeying the iden-
tity relation xn = 1, be bounded from above by a
certain natural number bed, n), depending only on d
and n? This is the so-called restricted Burnside problem.
It has been positively solved [13) for all prime
exponents n =p. It has thus been proved that there
exists a universal finite p-group B (d, p) of order bed, p)
whose quotient groups are isomorphic to all other finite
p-groups with d generators and satisfying the relation
x P = 1. If B (d, p) is finite, one has the identity
B(d,p)=B(d,p). A comparison of the results in [10)
and [13) leads to the conclusion that if p is sufficiently
large, there exists a finitely-generated infinite simple p-
group of exponent p. It has been shown that
b (2, 5) = 534 . For p ~ 7 there are only a few estimates
from below for b (d, p), connected with corresponding
estimates c(d, p) for the nilpotency class of the group
B (d, p). It is known that c(2, p) cannot be a linear
function of p. More importantly, c(d, p) increases
without limit with d [14), [15). The question of the
existence of B(d, n) for n=pm, m>I, beginning from
n = 8 and 9 is still (1977) open. At the same time, the
existence of B (d, n) for all square-free n is a conse-
quence of the results reported in [6) and [l3), and of
the theorem of the solvability of groups of odd-order
(cf. Burnside problem for finite groups), and of certain
facts concerning the classification of simple groups.
The original solutions [8) of the unbounded Burnside
problem and [l3) of the restricted Burnside problem are
partly based on the theory of algebras - in the former
case on the criterion of infinite dimensionality of an
algebra, in the latter on an identity of Lie algebras (cf.
Lie algebra), which is an analogue of the identity x P = 1
in groups [16), [17). There are other Burnside-type
problems than those mentioned above, and they receive
considerable attention [8), [9).
References
[I] BURNSIDE, W. ':)n an unsettled question in the theory of
discontinuous groups', Quart. J. Pure Appl. Math. 33 (1902),
230-238.
[2] LEVI, F.V. and W AERDEN, B.L. VAN DER: 'Ueber eine beson-
dere klasse von Gruppen', Abh. Math. Sem. Univ. Hamburg 9
(1932), 154-158.
[3] SANOV, LN.: 'Solution of the Burnside problem for exponent
4', Uchen. Zap. Leningrad. Gos. Univ. Ser. Mat. 10 (\ 940),
166-\70 (in Russian).
[4] BAYES, A.J., KAUTSKY, J. and WAMSLEY, T.W.: Proc. 2-nd
Internat. Con! Theory of Groups, Canberra, 1973, pp. 82-89.
[5] HALL, M.: 'Solution of the Burnside problem for exponent 6',
Proc. Nat. Acad. Sci. USA 43 (1957),751-753.
[6] HALL, P. and HIGMAN, G.: 'On the p-Iength of p soluble
groups and reduction theorems for Burnside's problem', Proc.
London Math. Soc. (3) 6 (1956), 1-42.
[7] NOVIKOV, P.S.: 'On periodic groups', A mer. Math. Soc. Transl.
487
BURNSIDE PROBLEM
45 (1965), 19-22 (in Russian). (Dokl. Akad. Nauk SSSR 127,
no. 4 (1959), 749-752)
(8) GoLOD, E.S.: 'On nil algebras and residually-finite groups',
Izv. Akad. Nauk SSSR Ser. Mat. 28 (1964),273-276 (in Rus-
sian).
(9) ALESHIN, S.V.: 'Finite automata and Burnside's problem for
periodic groups', Math. Notes 11, no. 3 (1972), 199-203. (Mat.
Zametki 11, no. 3 (1972), 319-328)
[lOA) NOVIKOV, P.S. and ADYAN, S.I.: Izv. Akad. Nauk SSSR Ser.
Mat. 32, no. I (1968),212-244.
[lOB) NOVIKOV, P.S. and ADYAN, S.I.: Izv. Akad. Nauk SSSR Ser.
Mat. 32, no. 3 (1968), 709-731.
[II) ADYAN, S.I.: The Burnside problem and identities in groups,
Springer, 1979 (translated from the Russian).
(12) SHMIDT, O.Yu.: Selected works on mathematics, Moscow, 1959,
pp. 298-300 (in Russian).
(13) KOSTRIKIN, A.I.: 'On Burnside's problem', Izv. Akad. Nauk
SSSR Ser. Mat. 23 (1959),3-34 (in Russian).
(14) BACHMUTH, S.M., MUCHIZUKI, H.Y. and WALKUP, D.W.: 'A
nonsolvable group of exponent 5', Bull. Amer. Math. Soc. 76,
no. 3 (1970), 638-640.
(15) RAzMYSLOV, Yu.P.: 'On Lie algebras satisfying the Engel con-
dition', Algebra and Logic 10, no. I (1971),21-29
(Algebra i Logika 10, no. I (1971),33-44)
(16) MAGNUS, W.: 'Ueber Beziehungen zwischen hoheren Komrnu-
tatoren', J. Reine Angew. Math. 177 (1937), 105-115.
(17) HIGMAN, G.: 'Lie ring methods in the theory of finite nilpotent
groups', in Proc. Internat. Congress Mathematicians Edinburgh,
1960, Cambridge Dniv. Press, 1960, pp. 307-312.
(18) KUROSH, A.G.: 'Ring-theoretical problems connected with
Burnside's problem for periodic groups', Izv. Akad. Nauk
SSSR Ser. Mat. 5 (1941),233-240 (in Russian). German
abstract.
A.l. Kostrikin
Editorial comments. A positive solution of the restricted
Burnside problem for the case of exponent five was first
obtained in [.A.11]. In the period 1977 - 1986 a lot of
additional results on the Burnside problem for periodic
groups were obtained. R.1. Grigorchuk [A1] proposed a
most-simple construction of finitely-generated infinite rr
groups. One of its versions leads, in particular (cf. [A2]), to
the construction of groups of intermediate growth, i.e. nei-
ther polynomial, nor exponential growth; it leads, moreover,
to the construction of both such periodic groups as well as
of such torsion-free groups (solutions to Milnor's problem).
A simple and geometrically obvious version of the nega-
tive solution to the bounded Burnside problem for odd
numbers n > 1010 was given by AYu. Ol'shanskiT [A3]. Later
(cf. [A4]) he constructed for each sufficiently-large prime
number p an infinite rrgroup in which all normal subgroups
have order p (Tarski's monster). This is the strongest form
of a negative solution to Burnside's problem. Every1hing
that has been done on the restricted Burnside problem is
summarized in [A5] and the book [A6]. In [A6] a complete
list of references is given, as well as results of computer
experiments of various scientists:
18(4,4) 1=2422 (cf. [A7]) ,
I 8(3,5) 1,,;;;52282 (cf. [A8]) ,
I 8(2, 7) I > 76366 ,
I 8(2,5) 1=534 (cf. [A9]).
Further, [A10] may prove to be useful.
References
[A 1) GRIGORCHUK, R.I.: 'On Burnside's problem for periodic
groups', Funktsional. Anal. i Prilozhen. 14, no. 1 (1980), 53-
54 (in Russian).
[A2) GRIGORCHUK, R.I.: 'Milnor's problem on the growth of
groups', Ookl. Akad. Nauk SSSR 271, no. 1 (1983), 30-33
(in Russian).
[A3] OL'SHANSKIT, A.Yu.: 'On a theorem of Novikov-Adyan', Mat.
Sb. 118, no. 2 (1982), 202-235 (in Russian).
[A4] OL'SHANSKIT, A.Yu.: 'Groups of bounded period in which all
subgroups are of prime order', Algebra i Logika 21, no. 5
(1982), 555-618 (in Russian).
[A5] VAUGHAN-LEE, M.R.: 'The restricted Burnside problem', Bull.
London Math. Soc. 17 (1985), 113-133.
[A6] KOSTRIKIN, A.I.: Around Burnside, Springer. Forthcoming.
[A7] ALFORD, HAVAS, G. and NEWMAN, M.F.: 'Groups of
exponent four', Notices Amer. Math. Soc. 22 (1975), A.301.
[A8] HAVAS, G., NEWMAN, M.F. and VAUGHAN-LEE, M.R.: 'A nil-
potent quotient algorithm for graded Lie rings', To appear.
[A9] HAVAS, Go, WALL, GoE. and WAMSLEY, loW.: 'The two-
generator restricted Burnside group of order five', Bull. Aus-
tral. [,lath. Soc. 10 (1974), 459-470.
[A 10] MENNICKE, l.L (ED.): 'Proc. Burnside workshop Bielefeld,
1977', Lec!. Notes Math., Vol. 806, Springer, 1980.
[A 11] HIGMAN, G.: 'On finite groups of exponent five', Proc. Cam-
bridge Phil. Soc 52 (1956), 381-390.
AMS 1980 Subject Classification: 20F50